Lecture notes on Planetary sciences and Satellite Orbit

Determination

Ernst J.O. Schrama

Delft University of Technology,
Faculty of Aerospace, Astrodynamics and Satellite missions
e-mail: e.j.o.schrama@tudelft.nl

29-Aug-2017
Contents

1 Introduction 7

2 Two body problem 10

2.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 10
2.2 Keplerian model . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11
2.2.1 Equations of motion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11
2.2.2 Keplerian equations of motion . . . . . . . . . . . . . . . . . . . . . . . . 13
2.2.3 Orbit plane . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 16
2.2.4 Parabolic and hyperbolic orbits . . . . . . . . . . . . . . . . . . . . . . . . 19
2.2.5 The vis-viva equation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
2.2.6 Keplers orbit in three dimensions . . . . . . . . . . . . . . . . . . . . . . 22
2.3 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 23

3 Potential theory 25
3.1 Solutions of the Laplace equation . . . . . . . . . . . . . . . . . . . . . . . . . . . 26
3.2 Legendre Functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 27
3.3 Normalization . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 29
3.4 Properties of Legendre functions . . . . . . . . . . . . . . . . . . . . . . . . . . . 29
3.4.1 Property 1 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 29
3.4.2 Property 2 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 30
3.4.3 Property 3 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 30
3.5 Convolution integrals on the sphere . . . . . . . . . . . . . . . . . . . . . . . . . . 31
3.6 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 32

4 Fourier frequency analysis 33

4.1 Continuous Fourier Transform . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 33
4.2 Discrete Fourier Transform . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 36
4.2.1 Fast Fourier Transform . . . . . . . . . . . . . . . . . . . . . . . . . . . . 36
4.2.2 Nyquist theorem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 37
4.2.3 Convolution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 37
4.2.4 Effect of a data selection window . . . . . . . . . . . . . . . . . . . . . . . 40
4.2.5 Parseval theorem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
4.3 Demonstration in MATLAB . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
4.3.1 FFT of a test function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
4.3.2 Harmonics of a sawtooth function . . . . . . . . . . . . . . . . . . . . . . 44
4.3.3 Gibbs effect . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 46

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 4.4 Power density spectra . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 47
4.5 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 48

5 Reference Systems 49
5.1 Definitions of ECI and ECF . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 50
5.1.1 Transformations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 51
5.1.2 Role of the IERS . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 52
5.1.3 True of date systemen . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 52
5.2 Representation problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 53
5.2.1 Geocentric coordinates . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 53
5.2.2 Geodetic coordinates . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 53
5.2.3 Map coordinates . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
5.3 What physics should we use? . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
5.3.1 Relativistic time dilatation . . . . . . . . . . . . . . . . . . . . . . . . . . 58
5.3.2 Gravitational redshift . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 60
5.3.3 Schwarzschild en Lense-Thirring on satellite orbit dynamics . . . . . . . . 62
5.4 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 63

6 Observation and Application 64

6.1 Techniques . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 64
6.1.1 Satellite Laser Ranging . . . . . . . . . . . . . . . . . . . . . . . . . . . . 66
6.1.2 Doppler tracking . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 66
6.1.3 Global Positioning System . . . . . . . . . . . . . . . . . . . . . . . . . . . 67
6.2 Corrections . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 72
6.2.1 Light time effect . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 72
6.2.2 Refraction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 72
6.2.3 Multipath . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 76
6.3 Applications . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 77
6.3.1 Satellite altimetry . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 77
6.3.2 Very long baseline interferometry . . . . . . . . . . . . . . . . . . . . . . . 83
6.3.3 Satellite gravimetry . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 86
6.4 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 89

7 Observations, models and parameters 90

7.1 Dynamical systems and statistics . . . . . . . . . . . . . . . . . . . . . . . . . . . 91

8 Parameter estimation 93
8.1 Random variables . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 93
8.1.1 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 93
8.1.2 Conditional probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . 94
8.1.3 Bayes theorem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 95
8.2 Probability density function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 97
8.3 Expectation operator . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 98
8.4 Covariance analysis . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 99
8.4.1 Covariance matrices in more dimensions . . . . . . . . . . . . . . . . . . . 99
8.4.2 Empirical orthogonal functions . . . . . . . . . . . . . . . . . . . . . . . . 102

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 8.4.3 Transformation of covariance matrices . . . . . . . . . . . . . . . . . . . . 103
8.5 Least squares method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 103
8.5.1 Parameter covariance matrix . . . . . . . . . . . . . . . . . . . . . . . . . 104
8.5.2 Example problem 1 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 105
8.6 Non-linear parameter estimation . . . . . . . . . . . . . . . . . . . . . . . . . . . 105
8.6.1 Example problem 2 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 106
8.7 Properties of the least squares algorithm . . . . . . . . . . . . . . . . . . . . . . . 108
8.7.1 Effect of scaling . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 108
8.7.2 Penrose-Moore pseudo-inverse . . . . . . . . . . . . . . . . . . . . . . . . . 108
8.7.3 Application of singular value decomposition . . . . . . . . . . . . . . . . . 109
8.8 Advanced topics . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 109
8.8.1 Compatibility conditions of Ax = y . . . . . . . . . . . . . . . . . . . . . . 109
8.8.2 Compatibility conditions At Ax = At y . . . . . . . . . . . . . . . . . . . . 110
8.8.3 Example problem 3 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 111
8.8.4 Constraint equations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 112
8.8.5 The Levenberg-Marquardt method . . . . . . . . . . . . . . . . . . . . . . 113
8.9 Implementation of algorithms . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 114
8.9.1 Accumulation algorithm, solve on the fly . . . . . . . . . . . . . . . . . . . 115
8.9.2 Sparse matrix solvers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 115
8.9.3 Blocking of parameters . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 115
8.9.4 Iterative equation solvers . . . . . . . . . . . . . . . . . . . . . . . . . . . 116

9 Modeling two-dimensional functions and data with polynomials 119

9.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 119
9.2 Polynomials to fit data and functions . . . . . . . . . . . . . . . . . . . . . . . . . 120
9.3 Modeling Doppler data . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 121
9.4 Fitting continuous and differentiable functions . . . . . . . . . . . . . . . . . . . . 122
9.5 Example continuous function fit . . . . . . . . . . . . . . . . . . . . . . . . . . . . 125
9.6 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 126

10 Dynamical parameter estimation 127

10.1 Laplace transformation approach . . . . . . . . . . . . . . . . . . . . . . . . . . . 127
10.1.1 Laplace Transforms demonstrated . . . . . . . . . . . . . . . . . . . . . . 129
10.1.2 Define and optimize the parameters . . . . . . . . . . . . . . . . . . . . . 131
10.2 Shooting problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 131
10.2.1 Numerical integration methods . . . . . . . . . . . . . . . . . . . . . . . . 133
10.2.2 Variational equations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 137
10.3 Parameter estimation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 137
10.3.1 The difference between dart and curling . . . . . . . . . . . . . . . . . . . 137
10.3.2 Numerical methods applied to the variational equations . . . . . . . . . . 139
10.3.3 Partial derivatives with respect to the control parameters . . . . . . . . . 140
10.3.4 Batch parameter estimation . . . . . . . . . . . . . . . . . . . . . . . . . . 141
10.3.5 Sequential parameter estimation . . . . . . . . . . . . . . . . . . . . . . . 143
10.3.6 Toy Kalman filter without process noise . . . . . . . . . . . . . . . . . . . 145
10.3.7 Toy Kalman filter with process noise . . . . . . . . . . . . . . . . . . . . . 147
10.3.8 Characterization of process noise . . . . . . . . . . . . . . . . . . . . . . . 147

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 10.3.9 Allan variance analysis . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 148

11 Three body problem 151

11.1 The restricted three-body problem . . . . . . . . . . . . . . . . . . . . . . . . . . 151
11.2 Two bodies orbiting the Sun . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 152
11.3 Accelerations in a rotating system . . . . . . . . . . . . . . . . . . . . . . . . . . 153
11.4 Jacobi constant . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 155
11.5 Position Lagrange points . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 156
11.6 Stability conditions in Lagrange points . . . . . . . . . . . . . . . . . . . . . . . . 161
11.7 Exercise . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 161

12 Co-rotating orbit dynamics 163

12.1 Solution of the Hill equations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 165
12.1.1 Homogeneous part . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 165
12.1.2 Particular solution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 166
12.2 Characteristic solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 166
12.2.1 Homogeneous solution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 167
12.2.2 Particular solution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 167
12.2.3 Particular resonant solution . . . . . . . . . . . . . . . . . . . . . . . . . . 167
12.3 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 168

13 Hill sphere and Roche limit 169

13.1 Hill sphere . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 169
13.2 Roche limit . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 169
13.3 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 171

14 Tide generating force 172

14.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 172
14.2 Tide generating potential . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 172
14.2.1 Proof . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 173
14.2.2 Work integral . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 174
14.2.3 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 176
14.2.4 Some remarks . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 176
14.3 Frequency analysis of observed tides . . . . . . . . . . . . . . . . . . . . . . . . . 176
14.3.1 Darwin symbols and Doodson numbers . . . . . . . . . . . . . . . . . . . 177
14.3.2 Tidal harmonic coefficients . . . . . . . . . . . . . . . . . . . . . . . . . . 177
14.4 Tidal harmonics . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 179
14.5 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 180

15 Tides deforming the Earth 183

15.1 Solid Earth tides . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 183
15.2 Long period equilibrium tides in the ocean . . . . . . . . . . . . . . . . . . . . . . 184
15.3 Tidal accelerations at satellite altitude . . . . . . . . . . . . . . . . . . . . . . . . 185
15.4 Gravimetric solid earth tides . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 186
15.5 Reference system issues . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 187
15.6 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 188

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16 Ocean tides 189
16.1 Equations of motion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 191
16.1.1 Newtons law on a rotating sphere . . . . . . . . . . . . . . . . . . . . . . 191
16.1.2 Assembly step momentum equations . . . . . . . . . . . . . . . . . . . . . 192
16.1.3 Advection . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 195
16.1.4 Friction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 196
16.1.5 Turbulence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 196
16.2 Laplace Tidal Equations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 197
16.3 Helmholtz equation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 199
16.4 Drag laws . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 200
16.5 Linear and non-linear tides . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 200
16.6 Dispersion relation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 201
16.7 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 203

17 Data analysis methods 204

17.1 Harmonic Analysis methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 204
17.2 Response method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 206
17.3 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 206

18 Load tides 208

18.1 Green functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 208
18.2 Loading of a surface mass layer . . . . . . . . . . . . . . . . . . . . . . . . . . . . 208
18.3 Computing the load tide with spherical harmonic functions . . . . . . . . . . . . 210
18.4 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 211

19 Altimetry and tides 212

19.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 212
19.2 Aliasing . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 212
19.3 Separating ocean tide and load tides . . . . . . . . . . . . . . . . . . . . . . . . . 212
19.4 Results . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 213
19.5 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 214

20 Tidal Energy Dissipation 215

20.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 215
20.2 Tidal energetics from lunar laser ranging . . . . . . . . . . . . . . . . . . . . . . . 217
20.2.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 217
20.2.2 Relevant observations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 217
20.2.3 Modeling . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 217
20.2.4 Interpretation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 219
20.3 Tidal energetics and fluid dynamics . . . . . . . . . . . . . . . . . . . . . . . . . . 220
20.3.1 Dissipation terms in the Laplace Tidal equations . . . . . . . . . . . . . . 220
20.3.2 A different formulation of the energy equation . . . . . . . . . . . . . . . 221
20.3.3 Integration over a surface . . . . . . . . . . . . . . . . . . . . . . . . . . . 222
20.3.4 Global rate of energy dissipation . . . . . . . . . . . . . . . . . . . . . . . 222
20.4 Rate of energy dissipation obtained from ocean tide models . . . . . . . . . . . . 225
20.4.1 Models . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 225

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 20.4.2 Interpretation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 226
20.5 Local estimates of tidal energy dissipation, internal waves and mixing . . . . . . 227
20.6 Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 227

6
Chapter 1

Introduction

In these lecture notes I bundled all material that I use for the introductory and advanced course
on planetary sciences and the course on satellite orbit determination which are part of the
curriculum at the faculty of aerospace engineering at the Delft University of technology. In the
MSc track of the faculty of aerospace engineering the course code for planetary sciences I is
ae4-890 and for the follow up course it is ae4-876. In the same curriculum the course on satellite
orbit determination comes with the code ae4-872.
A main topic in satellite orbit determination is the problem of parameter estimation which
is relates the dynamics of a space vehicle to observation techniques. From this follows a number
of scientific applications that are related to the observation techniques. In order to set-up the
framework for all lectures we start in chapter 2 with the two body problem, this material is also
mandatory for the planetary science I (ae4-890). It depends on your prior education whether or
not you need to study this chapter. Chapter 2 contains all required information from the BSc
of aerospace engineering. If you dont feel familiar with the two-body problem then study it
entirely, if you want to test your knowledge then try the exercises at the end of this chapter.
The two-body problem is directly related to potential theory which is nowadays most likely
not part of your bachelor program; for this reason Ive included chapter 3. For course ae4-
890 I do recommend to study the Laplace equation and the series expansion of the function
1/r in Legendre functions. An advanced topic is that the Laplace equation also comes with
higher order expansions in the potential functions. A summary of some well known properties
of Legendre functions, spherical harmonics and convolution integrals on the sphere should be
seen as a reference, that is, you should recognize spherical harmonics, and potential coefficients,
but you are not asked to reproduce for instance recursive relations of the Legendre functions.
At various points in the lecture notes we refer to Fouriers method of frequency analysis which
is described in chapter 4, this is a general mathematical procedure of which the results are used
throughout the lecture notes. It finds its application in tidal theory, it relates for instance to
the chapter 3 on potential theory where we mentioned the convolution on the sphere, and the
solution of the Hill equations in chapter 12 depends on Fourier series which are a characteristic
solution of the system. During various lectures I noticed that Fouriers method for frequency
analysis is often not part of the BSc curriculum, so I added the topic to these lecture notes. We
treat the continuous case to introduce the topic, but rapidly switch to the discrete case which
seems most applicable to what most people use. I included a number of examples in MATLAB
to demonstrate various properties related to the Fourier transforms.
The definition of time and coordinates is essential for all lectures; this topic is not part of

7
the curriculum of aerospace and for this reason I added chapter 5. This chapter discusses the
relation between the Earth center inertial (ECI) and the Earth center fixed (ECF) frame, the
role of the International Earth Rotation Service (IERS), and transformations between reference
systems. Other topics in this chapter are map projections and the consequence of special and
general relativity on the definition of time and coordinates.
In chapter 6 we discuss observation techniques and applications relevant for ae4-872. We
introduce satellite laser ranging (SLR), Doppler tracking (best known is the French DORIS
system) and the Global Positioning System (GPS). There are a number of corrections common
to all observation techniques, for this reason we speak about the light time effect, but also
refraction in the atmosphere and the ionosphere and including the phenomenon multipath which
is best known during radio tracking. The applications that we discuss are satellite altimetry,
very long baseline interferometry (VLBI) and satellite gravimetry.
For the course on satellite orbit determination I recommend to study chapter 7 where we
introduce the concept of combining observations, models and parameters, the material presented
here continues with what was presented in chapters 2 to 6. In section 7.1 we discuss the need to
consider dynamics when we estimate parameters. This brings us to chapter 8 where parameter
estimation techniques are considered without consideration of a dynamical model. The need for a
statistical approach is introduced for instance in 8.1 where the expectancy operator is defined in
8.3. With this knowledge we can continue to the least squares methods for parameter estimation
as discussed in 8.5. Chapter 10 discusses dynamical systems, Laplace transformations to solve
the initial value problem, shooting problems to solve systems of ordinary differential equations,
dynamical parameter estimation, batch and sequential parameter estimation techniques, the
Kalman filter and process noise and Allan variance analysis.
For ae4-890 we recommend to study the three-body problem which is introduced in chap-
ter 11. Related to the three-body problem is the consideration of co-rotating coordinate frames
in orbital dynamics, in these notes you can find this information in chapter 12, for the course on
ae4-890 we need this topic to explain long periodic resonances in the solar system, but also to
explain the problem of a Hill sphere which is found in [11]. During the lectures on solar system
dynamics in ae4-890 the Hill sphere and the Roche limit will be discussed in chapter 13 Both
topics relate to the discussion in chapters 2 and 13 of the planetary sciences book, cf. [11].
Course ae4-890 introduces the tide generating force, the tide generating potential and global
tidal energy dissipation. I recommend to study chapter 14 where we introduce the concept
of a tide generating potential whose gradient is responsible for tidal accelerations causing the
solid Earth and the oceans to deform. For planetary sciences II (ae4-876) I recommend the
remaining chapters that follow chapter 14. Deformation of the entire Earth due to an elastic
response, also referred as solid Earth tides and related issues, is discussed in chapter 15. A good
approximation of the solid Earth tide response is obtained by an elastic deformation theory.
The consequence of this theory is that solid Earth tides are well described by equilibrium tides
multiplied by appropriate scaling constants in the form of Love numbers that are defined by
spherical harmonic degree.
In ae4-876 we discuss ocean tides that follow a different behavior than solid earth tides.
Hydrodynamic equations that describe the relation between forcing, currents and water levels
are discussed in chapter 16. This shows that the response of deep ocean tides is linear, meaning
that tidal motions in the deep ocean take place at frequencies that are astronomically determined,
but that the amplitudes and phases of the ocean tide follow from a convolution of an admittance
function and the tide generating potential. This is not anymore the case near the coast where

8
non-linear tides occur at overtones of tidal frequencies. Chapter 17 deals with two well known
data analysis techniques which are the harmonic analysis method and the response method for
determining amplitude and phase at selected tidal frequencies.
Chapter 18 introduces the theory of load tides, which are indirectly caused by ocean tides.
Load tides are a significant secondary effect where the lithosphere experiences motions at tidal
frequencies with amplitudes of the order of 5 to 50 mm. Mathematical modeling of load tides
is handled by a convolution on the sphere involving Green functions that in turn depend on
material properties of the lithosphere, and the distribution of ocean tides that rest on (i.e.
load) the lithosphere. Up to 1990 most global ocean tide models depended on hydrodynamical
modeling. The outcome of these models was tuned to obtain solutions that resemble tidal
constants observed at a few hundred points. A revolution was the availability of satellites
equipped with radar altimeters that enabled estimation of many more tidal constants. This
concept is explained in chapter 19 where it is shown that radar observations drastically improved
the accuracy of ocean tide models. One of the consequences is that new ocean tide models result
in a better understanding of tidal dissipation mechanisms.
Chapter 20 serves two purposes, the section on tidal energetics from lunar laser ranging is
introduced in ae4-890, all material in section 20.2 should be studied for ae4-890. The other
sections in this chapter belong to course ae4-876, they provide background information with
regard to tidal energy dissipation. The inferred dissipation estimates do provide hints on the
nature of the energy conversion process, for instance, whether the dissipations are related to
bottom friction or conversion of barotropic tides to internal tides which in turn cause mixing of
between the upper layers of the ocean and the abyssal ocean.
Finally, while writing these notes I assumed that the reader is familiar with mechanics,
analysis, linear algebra, and differential equations. For several exercises we use MATLAB or
an algebraic manipulation tool such as MAPLE. There are excellent primers for both tools,
mathworks has made a matlab primer available, cf. [37]. MAPLE is suitable mostly for analysis
problems and a primer can be found in [35]. Some of the exercises in these notes or assigned as
student projects expect that MATLAB and MAPLE will be used.

E. Schrama, Delft September 29, 2017

9
Chapter 2

Two body problem

2.1 Introduction
The first astronomic observations were made more than two millennia ago, the quality of the
observations was constrained to the optical resolution and the sensitivity of the human eye. The
brightness of a star is usually indicated by its magnitude, a change of 1 in magnitude corresponds
to a change 2.5 in brightness. Under ideal conditions the human eye is limited to magnitude
six, and the optical resolution is roughly 15 (thus 15/3600 of a degree), while the angular
resolution of binoculars is 2.5. The naked eye is already a very sensitive and high quality
optical instrument for basic astronomic observations, as long as there is no light pollution and
when your eyes are used to darkness. We are able to distinguish planets from Mercury to Saturn,
comets, meteors and satellites but our naked-eye lacks the resolution to observe the moons of
Jupiter, or the second star of Mizar in Ursa Major.
The discussion about the motion of planets along the night sky goes back to ancient history.
The Greeks and Romans associated the planets with various gods. Mars was for instance the God
of War, Jupiter held the same role as Zeus in the Greek Pantheon and Mercury was the God of
trade, profit and commerce. Planets are unique in the night sky since the wander relative to the
stars, who seem to be fixed on a celestial sphere for an observer on a non-rotating Earth. Before
the invention of the telescope in 1608 and its first application for astronomic observations in
1610 by Galileo Galilei the believe was essentially that the Earth was the center of the universe,
that it was flat and that you could fall over the horizon and that everything else in the universe
rotated around the Earth.

Galileo, Copernicus, Brahe and Kepler

Galileo Galilei was an Italian astronomer (1564 to 1642) renowned for his revolutionary new
concept the solar system causing him to get into trouble with the inquisition. He modified the
then existing telescope into an instrument suitable for astronomic observations to conclude in
1610 that there are four Moons orbiting the planet Jupiter. The telescope was earlier invented
by the German-born Dutch eyeglass maker Hans Lippershey who demonstrated the concept
of two refracting lenses to the Dutch parliament in 1608. After all it is not surprising that
the observation of moons around Jupiter was made in southern Europe, which on the average
has a higher chance of clear night skies compared to the Netherlands. One of Galileo Galileis
comments on the classical view on the solar system was that his instrument permitted him to

10
see moons orbiting another planet, and that the classical model was wrong.
Other developments took place around the same time in Europe. Nicolaus Copernicus was
a Polish astronomer who lived from 1473 to 1543 and he formulated the concept of planets
wandering in circular orbits about the Sun, which was new compared to the traditional geocentric
models of Claudius Ptolomaeus (87 to 150) and the earlier model of Hypparchus (190 to 120 BC).
It was the Danish astronomer Tycho Brahe (1546 to 1601) to conclude on basis of observations of
the planet Mars that there were deviations from the Copernican model of the solar system. The
observations of Tycho Brahe assisted the German mathematician, astronomer and astrologer
Johannes Kepler 1571 to 1630) to complete a more fundamental model that explains the motion
of planets in our solar system. The Keplerian model is still used today because it is sufficiently
accurate to provide short-term and first-order descriptions of planetary ephemerides in our solar
system and satellites orbiting the Earth.

Keplers laws
The mathematical and physical model of the solar system ican be summarized in three laws
postulated by Kepler. The first and the second law were published in Astronomia Nova in 1609,
the third law was published in Harmonices Mundi in 1619:

Law I: In our solar system, the Sun is in a focal point of an ellipse, and the planets move
in an orbital plane along this ellipse, see plate 2.1.

Law II: The ratio of an area swept by a planet relative to the time required is a constant,
see plate 2.2.

Law III: The square of the mean orbital motion times the cube of the largest circle con-
taining the ellipse is constant. Thus:

n2 a3 = G.M = (2.1)

The constant n is the mean motion in radians per second and a the semi-major axis in some
unit of length. In this equation G is the universal gravitational constant and M is the mass of
the Sun. (both in units that correspond to the left hand side).

2.2 Keplerian model

In this section we demonstrate the validity of the Keplerian model, essentially by returning to
the equations of motion inside which we substitute a suitable gradient of a potential function.
This will result in an expression that describes the radius of the planet that depends on its
position in orbit. After this point we will derive a similar expression for the scalar velocity in
relation to the radius, the latter is called the vis-viva equation.

2.2.1 Equations of motion

In an inertial coordinate system the equations of motion of a satellite are:
X i
x = V + f (2.2)
i

11
 Figure 2.1: Elliptical orbit of a planet around the sun in one of the focal points

Figure 2.2: Keplers equal area law: segment AB-Sun and segment CD-Sun span equal areas,
the motion of the planet between A and B takes as long as it would between C and D

12
 i
where x is an acceleration vector and V a so-called potential function and where the terms f
represent additional accelerations. An in-depth discussion on potential functions can be found
in chapter 3. At this point it is sufficient to assume that the equations of motion in (2.2) apply
for a planet orbiting the sun. Equation (2.2) is a second-order ordinary differential equation
explaining that a particle in a force field is accelerating along the local direction of gravity (which
is the gradient of V written as V = (V /x, V /y, V /z) in the model). The model allows
for additional accelerations which are usually much smaller than the gravitational effect.
A falling object on Earth like a bullet leaving a gun barrel will exactly obey these equations.
In this case gravity is the main force that determines the motion, while also air drag plays a
significant role. One way to obtain a satellite in orbit would be to shoot the bullet with sufficient
horizontal velocity over the horizon. If there wouldnt be air drag then Keplers orbit model
predicts that this particular bullet eventually hits the gunman in his back. There are at least
two reasons why this will never happen. The first reason is of course the presence of air drag,
the second reason is that the coordinate frame we live in experiences a diurnal motion caused by
a rotation Earth. (It is up to you to verify that Keplers bullet will hit an innocent bystander
roughly 2000 km west of your current location on the equator.) Air drag will keep the average
bullet exiting a barrel within about 2 kilometer which is easy to verify when you implement
eq. (2.2) as a system of first-order ordinary differential equations in MATLAB. The rotating
Earth causes a much smaller effect and you will not easily notice it. (In reality cross-wind has
a more significant effect).
Foucaults pendulum is best used to demonstrate the consequences of rotating Earth. Jean
Bernard Leon Foucault was a French physicist who lived from 1819 to 1868 and he demonstrated
the effect of Earth rotation on a pendulum mounted in the Meridian Room of the Paris obser-
vatory in 1851, today the pendulum can be found in the Pantheon in Paris where it is a 28-kg
metal ball suspended by wire in the dome of this building. Foucaults pendulum will oscillate in
an orbital plane, due to the Coriolis forces that act on the pendulum we observe a steady shift
of this orbital plane that depends on the latitude of the pendulum. Some facts are:

The coordinate system used in equation (2.2) is an inertial coordinate system that does
not allow frame accelerations due to linear acceleration or rotation of the frame.

Whenever we speak about gravity on the Earths surface, as we all know it, we refer to the
sum of gravitational and rotational acceleration. Just gravitation refers to the acceleration
caused by Newtons gravity law.

The potential V in equation (2.2) is thus best referred to as a gravitational potential,

sometimes it is also called the geo-potential.

The concept of potential functions is best explained in a separate lecture on potential theory.
Chapter 3 describes some basic properties to arrive at a suitable potential function for the Kepler
problem.

2.2.2 Keplerian equations of motion

A suitable potential V for the Kepler model is:

V (r) = (2.3)
r

13
It is up to the reader to confirm that this function fulfills the Laplace equation, but also, that it
attains a value of zero at r = where r is the distance to the point mass and where = G.M
with G representing the universal gravitational constant and M the mass which are both positive
constants.
The gradient of V is the gravitational acceleration vector that we will substitute in the
general equations of motion (2.2), which in turn explains that a satellite or planet at (x, y, z)
will experience an acceleration (x, y, z) which agrees with the direction indicated by the negative
gradient of the potential function V = /r. The equations of motion in (2.2) may now be
rearranged as:
V X
x = + fxi
x
i
V X
y = + fyi (2.4)
y
i
V X
z = + fzi
z
i

which becomes:
x/t = x/r3 x/t = x
y/t = y/r3 y/t = y (2.5)
z/t = z/r3 z/t = z
In this case we have assumed that the center of mass of the system coincides with the origin. In
the three-body problem we will drop this assumption.

Demonstration of the gun bullet problem in matlab

In matlab you can easily solve equations of motion with the ode45 routine. This routine will
solve a first-order differential equation s = F (t, s) where s is a state vector. For a two body
problem we only need to solve the equations of motion in a two dimensions which are the in-plane
coordinates of the orbit. For the gun bullet problem we can assume a local coordinate system,
the x-axis runs away from the shooter and the y-axis goes vertically. The gravity acceleration
is constant, simply g = 9.81 m/ss . The state vector is therefore s = (x, y) and the gradient is
in this case V = (0, g) where g is a constant. In matlab you need to define a function to
compute the derivatives of the state vector, and in the command window you to call the ode45
procedure. Finally you plot your results. For this example we stored the function in a separate
file called dynamics.m containing the following code:

function [dsdt] = dynamics(t,s)

%
% in the function we will compute the derivatives of vector s
% with respect to time, the ode45 routine will call the function
% frequently when it solves the equations of motion. We store
% x in s(1) and y in s(2), and the derivatives go in s(3) and
% s(4). In the end dsdt receives the components of the
% gradient of V, here just (0,g)
%

14
dsdt = zeros(4,1); % we need to return a column vector to ode45
g = 9.81; % local gravity acceleration
dsdt(1) = s(3); % the velocity in the x direction is stored in s(3))
dsdt(2) = s(4); % the velocity in the y direction is stored in s(4))
dsdt(3) = 0; % there is no acceleration in the x direction
dsdt(4) = -g; % in the vertical direction we experience gravity

To invoke the integration procedure you should write another script that contains:

vel = 100; angle = 45;

s = [0 0 vel*cos(angle/180*pi) vel*cos(angle/180*pi)];
options = odeset(AbsTol,1e-10,RelTol,1e-10);
[T,Y] = ode45(@dynamics,[0 14],s,options );
plot(Y(:,1),Y(:,2))

The command s = ... assigns the initial state vector to the gun bullet, the options command is
a technicality, ie. probably you dont need it but when we model more complicated problems
then it may be needed. The odeset routine controls the integrator behavior. The next line calls
the integrator, and he last command plots the flight path of the bullet that we modelled. It
starts with a velocity of 100 m/s and the gun was aimed at 45 degrees into the sky, after about
14 seconds the bullet hits the surface 1000 meter away from the gun. Note that we did not
model any drag or wind effects on the bullet. In essence, all orbit integration procedures can be

Figure 2.3: Path of the bullet modelled in the script dynamics.m

15
Figure 2.4: The angular momentum vector is obtained by the cross product of the position and
velocity vector.

treated as variations of this problem, except that the dimension of the state vector will change
and that also, that the dynamics.m file will become more complicated.

2.2.3 Orbit plane

So far we have assumed that x y and z are inertial coordinates, and that the motion of the
satellite or planet takes place in a three dimensional space. The remarkable observation of
Kepler was that the motion occurs within a plane that intersects the center of the point source
mass generating V . This plane is called the orbit plane, and the interested reader may ask why
this is the case. To understand this problem we need to consider the angular momentum vector
H which is obtained as:
r v = x x = H (2.6)
where v is the velocity vector and r the position vector, see also figure 2.4. If we assume that
x = r = (x, y, 0) and that x = v = (x, y, 0) then:

x x 0
y y = 0
0 0 xy y x

which explains that the angular momentum vector is perpendicular to the plane spanned by r
and v. To demonstrate that H = 0 we evaluate:

x x = x x + x x
t
The last term is zero, due to the fact that:

x = x
r3

16
we also find that:
x x = 0

so that H = 0. A direct consequence is that we conserve angular momentum, and as we will
show later, we also conserve energy. The fact that the angular momentum vector is constant
in size and direction also explains why Kepler found an equal area law and that the motion is
confined to an orbital plane.

Substitution 1
To simplify the search for a solution we confine ourself to an orbital place. A convenient choice
is in this case to work in polar coordinates so that:

x = r cos
y = r sin

In the sequel we will substitute this expression in the equations of motion that follow from the
point mass potential, see also equation (2.5). An intermediate step is:

x = r cos r sin
y = r sin + r cos

so that:

x = r cos 2r sin r sin r2 cos

y = r sin + 2r cos + r cos r2 sin

which is equivalent to:

r r2
    
x cos sin
= (2.7)
y sin cos 2r + r
For the gradient we have:
    
V /x r/x /x V /r
= (2.8)
V /y r/y /y V /

so that:
cos sinr /r2
    
V /x
= cos (2.9)
V /y sin r
0
Since the right hand sides of (2.8) and (2.9) are equal we get:

r r2 = (2.10)
r2
2r + r = 0 (2.11)

For the length of the angular momentum vector we get:

h = |H| = xy y x
= +r cos (r sin + r cos ) r sin (r cos r sin )
= r2

17
which demonstrates that equal areas are covered in equal units of time in Keplers second law.
Since h is constant we obtain after differentiation with respect to time:

h = 2rr + r2 = 0 (2.12)

Since r = 0 is a trivial solution we keep:

2r + r = 0 (2.13)

which is equal to (2.11). This consideration does not lead to a new insight in the problem.
And thus we turn our attention to eq. (2.10) which we can solve with a new substitution of
parameters.

Substitution 2
At this point a suitable parameter substitution is r = 1/u and some convenient partial derivatives
are:
u 1
= 2
r r
u u r t 1 1 r2 r
= = ( 2 )(r)(1 ) = ( 2 )(r)( ) =
r t r r h h
2u u t r r r2 r
2
= ( ) = 1 = = 2 2
t h hh u h
from which we obtain:
2u
r = u2 h2 2

Substitution of these partial derivatives in (2.10) results in:

2 u h2
u2 h2 3 = u2
2 r
so that:
2u
2
+u= 2 (2.14)
h
This equation is equivalent to that of a mathematical pendulum, its solution is:

u = A cos + B
u
= A sin

2u
= A cos
2
We find:
2u
u+ 2
=B= 2
h
so that A becomes an arbitrary integration constant. In most textbooks we find the following
expression that relates r to :
a(1 e2 )
r() = (2.15)
1 + e cos

18
This expression results in circular orbits for e = 0, or elliptical orbits for 0 < e < 1. To verify
eq. (2.15) we evaluate r at the apo-apsis and the peri-apsis.
1
u( = 0) = = +A + B
a(1 e)
1
u( = ) = = A + B
a(1 + e)

From which we get:

e
A =
a(1 e2 )

B =
h2
1 1 2
2B = + =
a(1 e) a(1 + e) a(1 e2 )
1
B= 2
= 2
a(1 e ) h
resulting in: p
h= a(1 e2 )
which provides us with the length of the angular momentum vector.

2.2.4 Parabolic and hyperbolic orbits

So far we have demonstrated that circular and elliptic orbits appear, but in textbooks you also
find that parabolic and hyperbolic orbits exist as a solution of the Kepler problem. A parabolic
orbit corresponds to e = 1, and in a hyperbolic orbit e > 1. The parabolic orbit is one where
we arrive with a total energy of zero at infinity, therefore it is also called the minimum escape
orbit. Another option to escape the planet is to fly in a hyperbolic orbit, in this case we arrive
with a positive total energy at infinity. The total energy for the circular and eccentric Kepler
orbit is negative.

2.2.5 The vis-viva equation

Equation (2.15) contains all information to confirm Keplers first and second law. We will
now switch to an energy consideration of the Keplerian motion. Because of the conservation
of momentum we can not allow that energy disappears over time. This agrees with what we
observe in astronomy; planets and moons do not disappear on a cosmologic time scale (which is
only true if we leave tidal dissipation out of the discussion). If we assume that the total energy
of the system is conserved then:
1 m
mv 2 = d
2 r
where m and v represent mass and scalar velocity and where d is constant. We eliminate the
mass term m by considering d = d /m so that:

v2
=d+
2 r

19
The question is now to find d, since this would give us a relation to connect the scalar velocity in
an orbit to the radius r. This is what we call the vis-viva equation or the path-speed equation.
At the peri-apsis and the apo-apsis the velocity vectors are perpendicular to r. The length
of the moment vector (h) is nothing more than the product of the peri-apsis height and the
corresponding scalar velocity vp . The same property holds at the apo-apsis so that:

a(1 e)vp = a(1 + e)va (2.16)

The energy balance at apo-apsis and peri-apsis is:

va2 = 2d + 2 = 2d + 2 (2.17)
ra a(1 + e)

vp2 = 2d + 2 = 2d + 2 (2.18)
rp a(1 e)

From equation (2.16) we get:

 2
1+e
vp2 = va2 (2.19)
1e
This equation is substituted in (2.18):
 2
1+e
va2 = 2d + 2 (2.20)
1e a(1 e)

From this last equation and (2.17) you find:

 2    
1e
va2 = 2d + 2 = 2d + 2 (2.21)
1+e a(1 e) a(1 + e)

so that:

d=
2a
As a result we find that the total energy in the Kepler problem becomes:

v2
= (2.22)
2 r 2a
so that the total energy by mass for an object in orbit around a planet is constrained to:

Etot = (2.23)
2a
The scalar velocity of the satellite follows from the so-called vis-viva (Latin: living force1 )
relation: s  
2 1
v=
r a
which is an important relation that allows you to compute v as a function of r for a semi-major
axis a and a solar mass .
1
wikipedia mentions that the vis-viva is a obsolete scientific theory that served as an elementary and limited
early formulation of the principle of conservation of energy

20
Orbital periods
For a circular orbit with e = 0 and r = a we find that:
r

v=
a
If v = na where n is a constant in radians per second then:
r

na = = n2 a3
a
This demonstrates Keplers third law. Orbital periods for any parameter e [0, 1] are denoted
by and follow from the relation:
s
2 a3
= = 2
n
The interested reader may ask why this is the case, why do we only need to calculate the orbital
period of a circular orbit and why is there no need for a separate proof for elliptical orbits.
The answer to this question is already hidden in the conservation of angular momentum, and
related to this, the equal area law of Kepler. In an elliptical orbit the area dA of a segment
spent in a small time interval dt is (due to the conservation of angular momentum) equal to
dA = 12 h. The area A within the ellipse is:
Z 2
1
A= r()2 d (2.24)
=0 2

To obtain the orbital period we fit small segments dA within A, and we get:
Z 2 Z 2
r()2 2a2
= A/dA = d = 1 d = (2.25)
=0 h =0 a
which is valid for a > 0 and 0 e < 1. This demonstrates the validity of Keplers 3rd law.

Time vs True anomaly, solving Keplers equation

Variable in equation (2.15) is called the true anomaly and it doesnt progress linearly in time.
In fact, this is already explained when we discussed Keplers equal area law. The problem is
now that you need to solve Keplers equation which relates the mean anomaly M to an eccentric
anomaly E which in turn is connected via a goniometric relation to the true anomaly . The
discussion is rather mathematical, but over the centuries various methods have been developed
to solve Keplers equation. Without any further proof we present here a two methods to convert
the true anomaly , into an epoch t relative to the last peri-apsis transit t0 . The algorithms
assume that:
The mean anomaly M is defined as M = n.(t t0 ) where n is the mean motion in radians
per second for the Kepler problem.
The eccentric anomaly E relates to M via a transcendental relation: M = E e sin E.

The goniometric relation tan = 1 e2 sin E/(cos E e) is used to complete the con-
version of E to .

21
Iterative approach
There is an iterative algorithm that starts with E = M as an initial guess. Next we evaluate
Ei = M e sin Ei1 repeatedly until the difference Ei e sin Ei M converges to zero. The
performance of this algorithm is usually satisfactory in the sense that we obtain convergence
within 20 steps. For a given eccentricity e one may make a table with conversion values to be
used for interpolation. Note however that the iterative method becomes slow and that it may
not easily converge for eccentricities greater than 0.6.

Bessel function series

There are alternative procedures which can be found on the Wolfram website, cf. [29]. One
example is the expansion in Bessel functions:
M = E e sin E (2.26)
N
X 2
E = M+ Jn (n.e) sin(n.M ) (2.27)
n
1

The convergence of this series is relatively easy to implement in MATLAB. First you define
M between 0 and 2, and you assume a value for e and N . Next we evaluate E with the
series expansion and substitute the answer for M back in the first expression to reconstruct the
M that you started with. The difference between the input M , and the reconstructed M is
then obtained as a standard deviation for this simulation, it is an indicator for the numerical
accuracy. Figure 2.5 shows the obtained rms values when we vary e and N in the simulation.
The conclusion is that it is difficult to obtain the desired level of 1016 with just a few terms,
a series of N = 20 Bessel functions is convergent for e up to approximately 0.4, and N = 50 is
convergent for e up to approximately 0.5. In most cases we face however low eccentricity orbits
where e < 0.05 in which case there is no need to raise N above 5 or 10 to obtain convergence.
The Jn (x) functions used in the above expression are known as Bessel functions of the first
kind which are characteristic solutions of the so-called Bessel differential equation for function
y(x):
d2 y dy
x2 2 + x + (x2 2 )y = 0 (2.28)
dx dx
The Jn (x) functions are obtained when we apply the Frobenius method to solve equation (2.28),
the functions can be obtained from the integral:
Z
1
Jn (x) = (cos(n x sin( ))d (2.29)
0
More properties of the Jn (x) function can be found on the Wolfram website, also, the Bessel
functions are usually part of a programming environment such as MATLAB, or can be found in
Fortran or C/C++ libraries. Bessel functions of the first kind are characteristic solutions of the
Laplace equation in cylindrical harmonics which finds its application for instance in describing
wave propagation in tubes.

2.2.6 Keplers orbit in three dimensions

To position a Kepler orbit in a three dimensional space we need three additional parameters for
the angular momentum vector H. The standard solution is to consider an inclination parameter

22
Figure 2.5: Convergence of the Bessel function expansion to approximate the eccentric anomaly
E from the input which is the mean anomaly M between 0 and 2. The vertical scale is
logarithmic, the plateau is the noise floor obtained with a 8 byte floating point processor.

I which is the angle between the positive z-axis of the Earth in a quasi-inertial reference system
and H. In addition we define the angle that provides the direction in the equatorial plane
of the intersection between the orbit plane and the positive inertial x-axis, is also called the
right ascension of the ascending node. The last Kepler parameter is called , which provides
the position in the orbital plane of the peri-apsis relative to the earlier mentioned intersection
line.
The choice of these parameters is slightly ambiguous, because you can easily represent the
same Keplerian orbit with different variables, as has been done by Delauney, Gauss and others.
In any case, it should always be possible to convert an inertial position and velocity in three
dimension to 6 equivalent orbit parameters.

2.3 Exercises
Test your own knowledge:

1. What is the orbital period of Jupiter at 5 astronomical units? (One astronomical unit is
the orbit radius of the Earth)

2. Plot r(), v() and the angle between r() and v() for [0, 2] and for e = 0.01 and
a = 10000 km for = 3.986 1014 m3 s2 .

3. For an elliptic orbit the total energy is negative, for a parabolic orbit the total energy
is zero, ie. it is the orbit that allows to escape from Earth to arrive with zero energy at

23
 infinity. How do you parameterize parabolic orbits, how do you show that they are a
solution of the Kepler problem? How does this relate to the escape velocity on Earth?

4. Make a perspective drawing of the Kepler ellipse in 3D and explain all involved variables.

5. Design a problem to plot ground tracks for an arbitrary Kepler orbit, assume a constant
Earth rotation speed at a sidereal rate.

6. Implement the equations of motion for the Kepler orbit in matlab and verify the numerical
solution of r and v against the analytical formulas.

7. Demonstrate in matlab that the total energy is conserved for the Kepler problem. Your
starting point is an integrated trajectory.

24
Chapter 3

Potential theory

Potential fields appear in many forms in physics; in the case of solar system dynamics in planetary
sciences we consider usually potential functions related to the gravitational effect of a planet
or a star. But in physics you may also speak about magnetic or electric fields that are also
potential fields. A potential function describes the potential energy of an object at some point
in a gravitational field of another mass, which is usually the Sun or the Earth.1 Potential
energy of that object depends on the location of the object, but when we talk about the concept
potential function we refer to the normalized potential energy of the object in question without
consideration of its own mass. The gradient of the potential function is equal to the acceleration
vector predicted by Newtons gravity law. Yet, in the case of Newton we would have to deal
with vectors, now we can use a scalar function which reduces the complexity of the problem.
We consider the problem where we are moving around in a gravitational force field.2 Potential
energy relates to the problem of being somewhere in a force field, whereby the field itself is caused
by the gravitational attraction of a mass source that is usually far larger than the object moving
around this source. The potential at the end of the path minus the potential at the beginning
of the path is equal to the number of Joules per kg that we need to put in the motion that takes
place in this gravitational force field. If you move away from the source mass you have to push
the object, so you spend energy. But instead, when you approach the source mass then all this
potential energy comes back again for free, and if you move along surfaces of equal potential
energy then no extra energy is required to move around. Force fields that possess this property
are said to be conservative force fields.
Mathematically speaking this means that the Laplacian of the potential V is zero, and thus
that 2 V = 0. To explain why this is the case we go back to the Gauss integral theorem. The
theorem states that: Z I
(, w) d = (w, n) d 0 (3.1)
0
Here is the shape of an arbitrary body and 0 its
surface. Furthermore n is an vector of
length 1 that is directed outwards on a surface element, while w is an arbitrary vector function.
If we take w as the gradient of the potential V , and if we stay outside all masses that generate
1
Potential is related to the Latin word potentia which was used to describe political influence, power of strength.
2
Gravitation is the effect caused by the mass of the Sun or a planet, gravity is the effect that you experience
on a rotating planet.

25
V then: Z I
(, V ) d = (V, n) d 0 (3.2)
0
In a conservative force field the right hand side of this integral relation will vanish for any
arbitrary choice of that does not overlap with the masses that generate V . If we take an
infinitesimal small volume then the left hand side becomes:
2V 2V 2V
2 V = V = + + =0 (3.3)
x2 y 2 z 2
This equation is known as the Laplace equation, potential functions V that fulfill the Laplace
equation are said to generate a conservative force field V . And within such a conservative
force field you can always loop around along closed curves without losing any energy. Non-
conservative force fields also exist, in this case the opposite would happen, namely that you lose
energy along a closed path.
In physics all electric, magnetic and gravitational field are conservative. Gravitation is
unique in the sense that it doesnt interact with electric and magnetic fields. The latter two
fields do interact, the most general interaction between E and B is described by the Maxwell
equations that permit Electro-Magnetic waves. Gravitation does not permit waves, at least, not
in Newtonian physics. The theory of general relativity does allow for gravity waves, although
these waves have not yet been detected. Other effects caused by general relativity such as the
peri-helium precession of the planet Mercury or the gravitational bending of light have been
demonstrated. The concept gravity wave is also used in non-relativistic physics, and for
instance in the solution of the Navier Stokes equations. In this case we call a surface wave in a
hydrodynamic model a gravity wave because gravity is the restoring force in the dynamics.

3.1 Solutions of the Laplace equation

A straightforward solution of V that fulfills the Laplace equation is the function V = /r
where r is the radius of an arbitrary point in space relative to a source point mass. Later we
will show that this point mass potential function applies to the Kepler problem.
The minus sign in front of the gradient operator in equation 2.2 depends on the convention
used for the geopotential function V . If we start at the Earths surface the potential would
attain a value Va , and at some height above the surface it would be Vb . The difference between
Vb Va should in this case be positive, because we had to spend a certain number of Joules per
kilogram to get from a to b, and this can only be the case is Vb is greater than Va . Once we
traveled from the Earths surface to infinity there is no more energy required to move around,
because we are outside the potentia of the Earth. Thus we must demand that V = 0 at infinity.
The V = /r potential function is one of the many possible solutions of the Laplace
equation. We call it the point mass potential function. There are higher order moments of the
potential function. In this case we use series of spherical harmonics which are base functions
consisting of Legendre polynomials multiplied times goniometric functions. For the moment this
problem is deferred until we need to refine variations in the gravitational field that differ from
the central force field.

26
3.2 Legendre Functions
Legendre functions appear when we solve the Laplace equation (U = 0) by means of the
method of separation of variables. Normally the Laplace equation is transformed in spherical
coordinates r, , (r: radius, : longitude : co-latitude); this problem can be found in section
10.8 in [67] where the following solutions are shown:
U (r, , ) = R(r)G(, ) (3.4)
with:
1
R(r) = c1 rn + c2 (3.5)
rn+1
and where c1 and c2 are integration constants. Solutions of G(, ) appear when we apply
separation of variables. This results in so-called surface harmonics; in [67] one finds:
G(, ) = [Anm cos(m) + Bnm cos(m)] Pnm (cos ) (3.6)
where also Anm and Bnm are integration constants. The Pnm (cos ) functions are called associ-
ated Legendre functions and the indices n and m are called degree and order. When m = 0 we
deal with zonal Legendre functions and for m = n we are dealing with sectorial Legendre func-
tions, all others are tesseral Legendre functions. The following table contains zonal Legendre
functions up to degree 5 whereby Pn (cos ) = Pn0 (cos ):
P0 (cos ) = 1
P1 (cos ) = cos
3 cos 2 + 1
P2 (cos ) =
4
5 cos 3 + 3 cos
P3 (cos ) =
8
35 cos 4 + 20 cos 2 + 9
P4 (cos ) =
64
63 cos 5 + 35 cos 3 + 30 cos
P5 (cos ) =
128

Associated Legendre functions are obtained by differentiation of the zonal Legendre functions:
dm Pn (t)
Pnm (t) = (1 t2 )m/2 (3.7)
dtm
so that you obtain:
P11 (cos ) = sin
P21 (cos ) = 3 sin cos
P22 (cos ) = 3 sin2
 
15 3
P31 (cos ) = sin cos2
2 2
P32 (cos ) = 15 sin2 cos
P32 (cos ) = 15 sin3

27
Legendre functions are orthogonal base functions in an L2 function space whereby the inner
product is defined as: Z 1
Pn0 (x)Pn (x) dx = 0 n0 =
6 n (3.8)
1
and Z 1
2
Pn0 (x)Pn (x) dx = n0 = n (3.9)
1 2n + 1
In fact, these integrals are definitions of an inner product of a function space whereby Pn (cos )
are the base functions. Due to orthogonality we can easily develop an arbitrary function f (x)
for x [1, 1] into a so-called Legendre function series:

X
f (x) = fn Pn (x) (3.10)
n=0

The question is to obtain the coefficients fn when f (x) is provided in the interval x [1, 1].
To demonstrate this procedure we integrate on the right and left hand side of eq. 3.10 as follows:
Z 1 Z 1X
f (x)Pn0 (x) dx = fn Pn (x)Pn0 (x) dx (3.11)
1 1 n=0

Due to the orthogonality relation of Legendre functions the right hand side integral reduces to
an answer that only exists for n = n0 :
Z 1
2
f (x)Pn (x) dx = fn (3.12)
1 2n +1
so that: Z 1
2n + 1
fn = f (x)Pn (x) dx (3.13)
2 1
This formalism may be expanded in two dimensions where we now introduce spherical harmonic
functions:  a=1
cos m
Ynma (, ) = P (cos ) (3.14)
sin m a=0 nm
which relate to associated Legendre functions. In turn spherical harmonic functions possess
orthogonal relations which become visible when we integrate on the sphere, that is:
Z Z
4(n + m)!
Ynma (, )Yn0 m0 a0 (, ) d = (3.15)
(2n + 1)(2 0m )(n m)!
but only when n = n0 and m = m0 and a = a0 . Spherical harmonic functions Ynma (, ) are
the base of a function space whereby integral (3.15) defines the inner product. We remark
that spherical harmonic functions form an orthogonal set of basis functions since the answer of
integral (3.15) depends on degree n and the order m.
In a similar fashion spherical harmonic functions allow to develop an arbitrary function over
the sphere in a spherical harmonic function series. Let this arbirary function be called f (, )
and set as goal to find the coefficients Cnma in the series:
X
X n X
1
f (, ) = Cnma Ynma (, ) (3.16)
n=0 m=0 a=0

28
This problem can be treated in the same way as for the zonal Legendre function problem, in
fact, it is a general approach that may be taken for the subset of functions that can be developed
in a series of orthogonal (or orthonomal) base functions. Thus:
Z Z Z Z X
X n X
1
Yn0 m0 a0 (, )f (, ) d = Yn0 m0 a0 (, ) Cnma Ynma (, ) d (3.17)
n=0 m=0 a=0

which is only relevant when n = n0 and m = m0 and a = a0 . So that:

Z Z
1
Cnma = Nnm Ynma (, )f (, ) d (3.18)

where
4(n + m)!
Nnm = (3.19)
(2n + 1)(2 0m )(n m)!

3.3 Normalization
Normalization of Legendre functions is a separate issue that follows from the fact that we are
dealing with an orthogonal set of functions. There are several ways to normalize Legendre
functions, one choice is to rewrite integral (3.15) into a normalized integral:
Z Z
1
Y nma (, )Y n0 m0 a0 (, ) d = 1 (3.20)
4

where we simply defined new normalized functions with an overbar which are now called the
normalized spherical harmonic functions. It is obvious that they rely on normalized associated
Legendre functions:

(n m)! 1/2
 
P nm (cos ) = (2n + 1)(2 0m ) Pnm (cos ) (3.21)
(n + m)!
The use of normalized associated Legendre functions results now in an orthonormal set of spher-
ical harmonic base functions as can be seen from the new definition of the inner product in
eq. (3.20). It is customary to use the normalized functions because of various reasons, a very
important numerical reason is that stable recursive schemes for normalized associated Legendre
functions exist whereas this is not necessarily the case for the unnormalized Legendre functions.
This problem is beyond the scope of these lecture notes, the reader must assume that there is
software to compute normalized associated Legendre functions up to high degree and order.

3.4 Properties of Legendre functions

3.4.1 Property 1
A well-known property that we often use in potential theory is the development of the function
1/r in a series of zonal Legendre functions. We need to be a bit more specific on this problem.
Assume that there are two vectors p and q and that their length is rp and rq respectively. If the
length of the vector p q is called rpq then:

1/2
rpq = rp2 + rq2 2rp rq cos (3.22)

29
for which it is known that:

1 X rp n

1
= Pn (cos ) (3.23)
rpq rq rq
n=0
where is the angle between p and q. This series is convergent when rp < rq . The proof for
this property is given in [52] and starts with a Taylor expansion of the test function:

1/2
rpq = rp 1 2su + s2 (3.24)

where s = rq /rp and u = cos . The binomial theorem, valid for |z| < 1 dictates that:

(1 z)1/2 = 0 + 1 z + 2 z 2 + ... (3.25)

where 0 = 1 and n = (1.3.5...(2n 1))/(2.4...(2n)). Hence if |2su s2 | < 1 then:

(1 2su + s2 )1/2 = 0 + 1 (2su s2 ) + 2 (2su s2 )2 + ... (3.26)

so that:
3 1
(1 2su + s2 )1/2 = 1 + us + (u2 )s2 + ...
2 3
= P0 (u) + sP1 (u) + s2 P2 (u) + ...

which completes the proof.

3.4.2 Property 2
The addition theorem for Legendre functions is:
1 X
Pn (cos ) = Y nma (p , p )Y nma (q , q ) (3.27)
2n + 1 ma

where p and p are the spherical coordinates of vector p and q and q the spherical coordinates
of vector q.

3.4.3 Property 3
The following recursive relations exist for zonal and associated Legendre functions:
n1 2n 1
Pn (t) = Pn2 (t) + tPn1 (t) (3.28)
n n
Pnn (cos ) = (2n 1) sin Pn1,n1 (cos ) (3.29)
Pn,n1 (cos ) = (2n 1) cos Pn1,n1 (cos ) (3.30)
(2n 1) (n + m 1)
Pnm (cos ) = cos Pn1,m (cos ) Pn2,m (cos ) (3.31)
nm nm
Pn,m (cos ) = 0 for m > n (3.32)
For differentiation the following recursive relations exist:
dPn (t)
(t2 1) = n (tPn (t) Pn1 (t)) (3.33)
dt

30
3.5 Convolution integrals on the sphere
Spherical harmonic function expansions are very convenient for the evaluation of the following
type of convolution integrals on the sphere:
Z
H(, ) = F (0 , 0 )G() d (3.34)

where d = sin d d and the spherical distance between , and 0 , 0 and the azimuth.
Functions F and G are written as:
X
X n X
1
F (, ) = Fnma Y nma (, ) (3.35)
n=0 m=0 a=0

where

Y nm,0 (, ) = cos(m)P nm (cos )

Y nm,1 (, ) = sin(m)P nm (cos )

and

X
G() = Gn Pn (cos ) (3.36)
n=0

which takes the shape of a so-called Greens function3 . It turns out that instead of numerically
computing the expensive surface integral in eq. (3.34) that it is easier to multiply the Gn and
Fnma coefficients:
X
X n X 1
H(, ) = Hnma Y nma (, ) (3.37)
n=0 m=0 a=0
where
4Gn
Hnma = Fnma (3.38)
2n + 1
For completeness we also demonstrate the validity of eq. (3.38). The addition theorem of Leg-
endre functions states that:
n
1 X
Pn (cos pq ) = P nm (cos p )P nm (cos q ) cos(m(p q )) (3.39)
2n + 1
m=0

which is equal to
n X
1
1 X
Pn (cos pq ) = Y nm (p , p )Y nm (q , q ) (3.40)
2n + 1
m=0 a=0

When this property is substituted in eq. (3.34) then:

Z (X )(
X Gn 0
)
H(, ) = Fnma Y nma (0 , 0 ) Y n0 m0 a0 (, )Y n0 m0 a0 (0 , 0 ) d (3.41)
nma 0 0 0
2n0 + 1
nma
3
George Green (1793-1841)

31
which is equal to:
Z
X Gn0 X
H(, ) = Y n0 m0 a0 (, ) Fnma Y nma (0 , 0 )Y n0 m0 a0 (0 , 0 ) d (3.42)
2n0 + 1
n0 m0 a0 nma

Due to orthogonality properties of normalized associated Legendre functions we get the desired
relation:
X 4Gn
H(, ) = Fnma Y nma (, ) (3.43)
nma
2n + 1
which completes our proof.

3.6 Exercises
1
1. Show that U = r is a solution of the Laplace equation U = 0

2. Show that the gravity potential of a solid sphere is the same as that of a hollow sphere
and a point mass

3. Demonstrate in matlab that eq. (3.23) rapidly converges when rq = f rp where f > 1.1
for randomly chosen values of and rp

4. Demonstrate in matlab that eqns. (3.14) are orthogonal over the sphere
 
1 x [0, 1]
5. Develop a method in matlab to express the Greens function f (x) =
P 0
as a series of Legendre functions f (x) = n an Pn (x).

32
Chapter 4

Fourier frequency analysis

Jean-Baptiste Joseph Fourier (17681830) was a French scientist who introduced a method
of frequency analysis where one could approximate an arbitrary function by a series of sine
and cosine expressions. He did not show that the series would always converge, the German
mathematician Dirichlet (1805-1859) later showed that there are certain restrictions of Fouriers
method, in reality these restrictions are usually not hindering the application of Fouriers method
in science and technology. Fouriers frequency analysis method assumes that we analyze a
function on a defined interval, Fourier made the crucial assumption that the function repeats
itself when we take the function beyond the nominal interval. For this reason we say that the
function to analyze with Fouriers method is periodic.
In the sequel we consider a signal v(t) that is defined in the time domain [0, T ] where T is the
length in seconds, periodicity implies that v(t + kT ) = v(t) where k is an arbitrary integer. For
k = 1 we see that the function v(t) simply repeats because v(t) = v(t + T ), we see the same on
the preceding interval because v(t) = v(t T ). Naturally one would imagine a one-dimensional
wave phenomenon like what we see in rivers, in the atmosphere, in electronic circuits, in tides,
and when light or radio waves propagate. This is what Fouriers method is often used for, the
frequency analysis reveals how processes repeat themselves in time, but also in place or maybe
along a different projection of variables. This information is crucial for understanding a physical
or man-made signal hidden in often noisy observations.
This chapter is not meant to replace a complete course on Fourier transforms and Signal
Processing, but instead we present a brief summary of the main elements relevant for our lectures.
If you have never dealt with Fouriers method then study both sections in this chaper, and test
your own knowledge by making a number of assignments at the end of this chapter. In case you
already attended lectures on the topic then keep this chapter as a reference. In the following
two sections we will deal with two cases, namely the continuous case where v(t) is an analytical
function on the interval [0, T ] and a discrete case where we have a number of samples of the
function v(t) within the interval [0, T ]. Fouriers original method should be applied to the
continuous method, for data analysis we are more inclined to apply the discrete Fourier method.

4.1 Continuous Fourier Transform

Let v(t) be defined on the interval t [0, T ] where we demand that v(t) has a finite number of
oscillations and where v(t) is continuous on the interval. Fourier proposed to develop v(t) in a

33
series:
N/2
X
v(t) = Ai cos i t + Bi sin i t (4.1)
i=0

where Ai and Bi denote the Euler coefficients in the series and where variable i is an angular
rate that follows from i = i where = 2 T . At this point one should notice that:

The frequency associated with T1 is 1 Hertz (Hz) when T is equal to 1 second. A record
length of T = 1000 seconds will therefore yield a frequency resolution of 1 milliHertz
because of the definition of equation (4.1).

Fouriers method may also be applied in for instance orbital dynamics where T is rescaled
to the orbital period, in this case we speak of frequencies in terms of orbital periods, and
hence the definition cycles per revolution or cpr. But other definitions of frequency are
also possible, for instance, cycles per day (cpd) or cycles per century (cpc).

When v(t) is continuous there are an infinite number of frequencies in the Fourier series.
However, all Euler coefficients that you find occur at multiples of the base frequency 1/T .

A consequence of the previous property is that the spectral resolution is only determined
by the record length during the analysis, the frequency resolution f is by definition 1/T .
The frequency resolution f should not be confused with sampling of the function v(t) on
t [0, T ]. Sampling is a different topic that we will deal in section 4.2 where the discrete
Fourier transform is introduced.

In order to calculate Ai and Bi in eq. (4.1) we exploit the so-called orthogonality properties of
sine and cosine functions. The orthogonality properties are defined on the interval [0, 2], later
on we will map the interval [0, T ] to the new interval [0, 2] which will be used from now on.
The transformation from [0, T ] or even [t0 , t0 + T ] to [0, 2] is not relevant for the method at
this point, but is will become important if we try to assign physical units to the outcome of the
result of the Fourier transform. This is a separate topic that we will discuss in section 4.4. The
problem is now to calculate Ai and Bi in eq. (4.1) for which we will make use of orthogonality
properties of sine and cosine expression. A first orthogonality property is:
Z 2
sin(mx) cos(nx) dx = 0 (4.2)
0

This relation is always true regardless of the value of n and m which are both integer whereas
x is real. The second orthogonality property is:

Z 2 0 : m 6= n
cos(mx) cos(nx) dx = : m=n>0 (4.3)
0
2 : m=n=0

and the third orthogonality property is:

Z 2 
: m=n>0
sin(mx) sin(nx) dx = (4.4)
0 0 : m 6= n, m = n = 0

34
The next step is to combine the three orthogonality properties with the Fourier series definition
in eq. (4.1). We do this by evaluating the integrals:
Z 2  
cos(mx)
v(x) dx (4.5)
0 sin(mx)

where we insert v(t) but now expanded as a Fourier series:

Z 2 X N/2  
cos(mx)
An cos(nx) + Bn sin(nx) dx (4.6)
0 sin(mx)
n=0

You can reverse the summation and the integral, the result is that many terms within this
integral disappear because of the orthogonality relations. The terms that remain result in the
following expressions: Z 2
1
A0 = v(x) dx, B0 = 0 (4.7)
2 0
1 2
Z
An = v(x) cos(nx) dx, n > 0 (4.8)
0
1 2
Z
Bn = v(x) sin(nx) dx, n > 0 (4.9)
0
The essence of Fouriers frequency analysis method can now be summarized:

The conversion of time domain to frequency domain goes via three integrals where we
compute An and Bn that appear in eq. (4.1). This conversion or transformation step is
called the Fourier transformation and it is only possible when v(x) exists on the interval
[0, 2]. Fourier series exist when there are a finite number of oscillations between [0, 2],
this means that a function like sin(1/x) could not be expanded. A second condition
imposed by Dirichlet is that there are a finite number of discontinuities. The reality in
most data analysis problems is that we hardly ever encounter the situation where the
Dirichlet conditions are not met.

When we speak about a spectrum we speak about the existence of the Euler coefficients
An and Bn .
Euler coefficients are often taken together in a complex number Zn = An +jBn
where j = 1. We prefer the use of j to avoid any possible confusing with electric
currents.

There is a subtle difference between the discrete Fourier transform and the continuous
transform discussed in this section. The discrete Fourier transform introduces a new
problem, namely that or the definition of sampling, it is discussed in section 4.2.

The famous theorem of Dirichlet reads according to [67]: If v(x) is a bounded and periodic
function which in any one period has at most a finite number of local maxima and minima and
a finite number of point of discontinuity, then the Fourier series of v(x) converges to v(x) at all
points where v(x) is continuous and converges to the average of the right- and left-hand limits
of v(x) at each point where v(x) is discontinuous.

35
 If the Dirichlet conditions are met then we are able to define integrals that relate f (t) in the
time domain and g() in the frequency domain:
Z
f (t) = g()ejt d (4.10)

Z
1
g() = f ( )ej d (4.11)
2
In both cases we deal with complex functions where at each spectral line two Euler coefficients
from the in-phase term An and the quadrature term Bn . The in-phase nomenclature originates
from the fact that you obtain the coefficient by integration with a cosine function which has a
phase of zero on an interval [0, 2] whereas a sine function has a phase of 90 . The amplitude
of each spectral line is obtained as the length of Zn = An + jBn , thus |Zn | whereas the phase
is the argument of the complex number when it is converted to a polar notation. The phase
definition only exists because it is taken relative to the start of the data analysis window, this
also means that the phase will change if we shift that window in time. It is up to the reader to
show how the resulting Euler coefficients are affected.

4.2 Discrete Fourier Transform

The continuous case introduced the theoretical foundation for what you normally deal with as
a scientist or engineer who collected a number of samples of the function v(tk ) where tk =
t0 + (k 1)t with k [0, N 1] and t > 0. The sampling interval is now called t. The length
of the data record is thus T = k.t, the first sample of v(t0 ) will start at the beginning of the
interval, and the last sample of the interval is at T t because v(t0 + T ) = v(t0 ).
When the first computers became available in the 60s equations (4.7), (4.8) and (4.9) where
coded as shown. Equation (4.7) asks to compute a bias term in the series, this is not a lot
of work, but equations (4.8) and (4.9) ask to compute products of sines and cosines times the
input function v(tk ) sampled on the interval [t0 , t0 + (N 1)t]. This is a lot of work because
the amount of effort is like 2N multiplications for both integrals times the number of integrals
that we can expect, which is the number the frequencies that can be extracted from the record
[t0 , t0 + (N 1)t]. Due to the Nyquist theorem the number of frequencies is N/2, and for each
integral there are N multiplications: the effort is of the order of N 2 operations.

4.2.1 Fast Fourier Transform

There are efficient computer programs (algorithms) that compute the Euler coefficients in less
time than the first versions of the Fourier analysis programs. Cooley and Tukey developed in
1966 a faster method to compute the Euler coefficients, they claim that the number of operations
is proportional to O(N log N ). Their algorithm is called the fast Fourier transform, or the FFT,
the first implementation required an input vector that had 2k elements, later versions allowed
other lengths of the input vector where the largest prime factor should not exceed a defined
limit. The FFT routine is available in many programming languages (or environments) such as
MATLAB. The FFT function assumes that we provide it a time vector on the input, on return
you get a vector with Euler coefficients obtained after the transformation which are stored as
complex numbers. The inverse routine works the other way around, it is called iFFT which

36
stands for the inverse fast Fourier transform. The implementation of the discrete transforms in
MATLAB follows the same definition that you find in many textbooks, for FFT it is:
N
X 1
Vk = vn e2jkn/N with k N and vn C and Vk C (4.12)
n=0

and for the iFFT it is:

N 1
1 X
vn = Vk e2jkn/N with n N and vn C and Vk C (4.13)
N
k=0

where vn is in the time domain while Vk is in the frequency domain, furthermore Eulers formula
is used: ejx = cos x + j sin x. Because of this implementation in MATLAB a conversion is
necessary between the output of the FFT stored in Vk to the Euler coefficients that we defined
in equations (4.1) (4.7) (4.8) and (4.9), this topic is worked out in sections 4.3.1 and 4.3.2 where
we investigate test functions.

4.2.2 Nyquist theorem

The Nyquist theorem (named after Harry Nyquist, 1889-1976, not to be confused with the
Shannon-Nyquist theorem) says that the number of frequencies that we can expect in a discretely
sampled record [t0 , t0 + (N 1)t] is never greater than N/2. Any attempt to compute integrals
(4.8) and (4.9) beyond the Nyquist frequency will result in a phenomenon that we call aliasing
or faltung (in German). In general, when the sampling rate 1/t is too low you will get an
aliased result as is illustrated in figure 4.1. Suppose that your input signal contains power
beyond the Nyquist frequency as a result of undersampling, the result is that this contribution
in the spectrum will fold back into the part of the spectrum that is below the Nyquist frequency.
Figure 4.2 shows how a spectrum is distorted because the input signal is undersampled. Due
to the Nyquist theorem there are no more than N/2 Euler coefficient pairs (Ai , Bi ) that belong
to a unique frequency i , see also eq. (4.1). The highest frequency is therefore N/2 times the
base frequency 1/T for a record that contains N samples. If we take a N that is too small then
the consequence may be that we undersample the signal, because the real spectrum of the
N
signal may contain power above the cutoff frequency 2T imposed by the way we sampled the
signal. Undersampling results in aliasing so that the computed spectrum will appear distorted.
Oversampling is never a problem, this is only helpful to avoid that aliasing will occur, however,
sometimes oversampling is simply not an option. In electronics we can usually oversample, but
in geophysics etc we can not always choose the sampling rate the way we would like it. Frequency
resolution is determined by the record length, short records have a poor frequency resolution,
longer records often dont.

4.2.3 Convolution
To convolve is not a verb you would easily use in daily English, according to the dictionary
it means to roll or coil together; entwine. When you google for convolved ropes you get to
see what you find in a harbor, stacks of rope rolled up in a fancy manner. In mathematics
convolution refers multiplication of two periodic functions where we allow one function to shift

37
Figure 4.1: Demonstration of the aliasing, suppose that the true signal was in blue, but that
we sample the blue signal at the red circles, any Fourier analysis procedure will now think that
the signal is recovered as the green function. Of course the green function is not the real signal,
instead we say that it is the aliased function. Remedies are, collect samples of v(tk ) at a higher
rate or, as is done in the case of tides, assume that you know the frequency of the blue function
so that the amplitude and phase of the green function can be used to recover the blue function.

38
Figure 4.2: Demonstration of the aliasing. The true spectrum of your signal is in red, the graph
displays the power at each frequency computed as Pi = (A2i + Bi2 ). The Nyquist frequency
is defined by the sampling rate of the input signal. Since aliasing results in folding the red
spectrum is folded back across the black dashed line which coincides with the Nyquist frequency.
The part that aliases back is the blue dashed graph left of the dashed black line, it adds to the
true spectrum which was red, so that the result will be the blue spectrum which is said to be
affected by aliasing. To summarize the situation, red is the real signal, but blue is what we
recover because our sampling rate was too low.

39
along another during the operation:
Z
h(t) = f ( )g(t )d (4.14)

we also say the h(t) is the result of the convolution of f (t) and g(t), the function f (t) would be
for instance a signal and g(t) could be a filter, so that h(t) is the filtered version of the signal.
The problem with direct convolution in the time domain is that the process is very slow, but
fortunately we can make use of one of the properties of the Fourier transform that greatly speeds
up the evaluation of the convolution integral.

F () = F F T (f (t))
G() = F F T (g(t))
H() = F () G()
h(t) = iF F T (H())

where is frequency and t time. Convolution is used to build, design and analyze filters in
digital communication, in physics convolution is often the result of a physical property between
two quantities. Since the advent of the FFT transform it has become possible to quickly carry
out convolutions with the help of a computer. In this sense FFT is the enabling technology for
digital signal processing.

4.2.4 Effect of a data selection window

During an analysis of a finite length data record we always deal somehow with the problem
convolution. Reason is that the length of the record itself acts like a box window that we impose
on a perhaps much longer data record. It was the choice of an observer to select a certain part
of the much longer record, and as such we could also affect the spectrum that we compute by
the choice of our window. So the spectrum that we get to see will be affected by convolution
of the box window being our selected data window. There are several ways one can handle this
problem:
Pre-process the data by removing a long term trend function from the input signal using a
least squares regression technique, see section 8.5 for a discussion. Geophysical data may
for instance show a slow drift or it may be irregularly spaced and if we would analyze a
record without taking care of that drift (or bias) term then just the presence of a drift term
would add the spectrum of the sawtooth function, for a discussion see section 4.3.2 where
we compute its spectrum in MATLAB. This is not what we want to see, so we first remove
the trend function from the data to retain a difference signal that we subject to the FFT
method. Interpolation and approximation may be a part of the story, these methods help
to get the data presented to the FFT method in such a shape that it becomes regularly
spaced and detrended, for a discussion see chapter 9.

The signal spectrum may be such that there is a fair amount of red signal. A spectrum is
said to be red if it contains, in analogy with optics, a significant amount of energy at the
lower frequencies compared to the rest of the spectrum. When you select a data record
then it may be such that the record itself is not a multiple of the length contained in the
red part of the spectrum. This leads to distortion of the red peaks in the spectrum, instead

40
Figure 4.3: Convolution: The signal in the top figure is convolved with the signal in the middle
figure and the result is presented in the bottom figure. Imagine that you slide the narrow block
function in the middle figure along the upper block function that you hold in place. At each
step you carry out the multiplication and the summation, and the result is stored in the lower
graph, this is done for all possible settings of the shift that we apply. When the middle block is
entirely positioned under the top block a value of one is plotted in the bottom graph, when the
middle block is not under the top block a zero is found, and when the middle block is partially
under the top block a partial result is found. Since the overlapping area is linearly increasing
with the applied shift we get to see linear flanks on the convolved function. Please notice that
we applied a scaling factor to redefine the maximum of the h(t) function as 1, in reality it isnt.

41
 of being a sharp peak the content of those peaks may smear to neighboring frequencies.
This is what we call spectral leakage. A possible remedy is to apply a window or tapering
function to the input data prior to computing the spectrum.

The choice of a taper function is a rather specific topic, tapering means that we multiply a
weighting function wn times the input data vn which results in vn0 that we subject (instead of
vn ) to the FFT method:

vn0 = wn .vn where n [0, N 1] and {wn , vn , vn0 } R and {n, N } N (4.15)

The result will be that the FFT(v 0 ) will improve in quality compared to the FFT(v), one aspect
that would be improved is spectral leakage. There are various window functions, the best known
general purpose taper is the Hamming function where:

wn = 0.54 0.46 cos(2n/N ), 0 n N (4.16)

MATLABs signal processing toolbox offers a variety of tapering functions, the topic is too
detailed to discuss here.

4.2.5 Parseval theorem

In section 4.2.3 we demonstrated that multiplication of Euler coefficients of two functions in
the frequency domain is equal to convolution in the time domain. Apply now convolution of
a function with itself at zero shift and you arrive at Parsevals identity, after (Marc-Antoine
Parseval 1755-1836) which says that the sum of the squares in the time domain is equal to
the sum of the squares in the frequency domain after we applied Fouriers transformation to a
record in the time domain, see section 4.2.5. The theorem is relevant in physics, it says that
the amount of energy stored in the time domain can never be different from the energy in the
frequency domain: X X
F 2 () = f 2 (t) (4.17)
i

where F is the Fourier transform of f .

4.3 Demonstration in MATLAB

4.3.1 FFT of a test function
In MATLAB we work with vectors and the set-up is such that one can easily perform matrix
vector type of operations, the FFT and the iFFT operator are implemented as such, they are
called fft() and ifft(). With FFT(f (x)) it does not precisely matter how the time in x is defined,
the easiest assumption is that there is a vector f in MATLAB and that we turn it into a vector g
via the FFT, the command would be g = fft(f ) where f is evaluated at x that appear regularly
spaced in the domain [0, 2], thus x = 0 : 2/N : 2 2/N in MATLAB. Before you blindly
rely on a FFT routine in a function library it is a good practice to subject it to a number of
tests. In this case we consider a test function of which the Euler coefficients are known:

f (x) = 7 + 2 sin(3x) + 4 cos(12x) 5 sin(13x); with x [0, 2] (4.18)

42
A Fourier transform of f should return to us the coefficients 7 at the zero frequency, 2 at the 3rd
harmonic, +4 at the 12th harmonic and -5 at the 13th harmonic. The term harmonic comes from
communications technology and its definition may differ by textbook, we say that the lowest
possible frequency at 1/T that corresponds to the record length T equals to the first harmonic,
at two times that frequency we have the second harmonic, and so on. I wrote the following
program in MATLAB to demonstrate the problem:
clear;
format short
dx = 2*pi/1000; x = 0:dx:2*pi-dx;
f = 2*sin(3*x) + 5 + 4*cos(12*x) - 5*sin(13*x);
g = fft(f);
idx = find(abs(g)>1e-10);
n = size(idx,2);
K = 1/size(x,2);
for i=1:n,
KK = K;
if (idx(i) > 1),
KK = 2*K;
end
A = KK*real(g(idx(i)));
B = KK*imag(g(idx(i)));
fprintf(%4d %12.4f %12.4f\n,[idx(i) A B]);
end
The output that was produced by this program is:
1 5.0000 0.0000
4 0.0000 -2.0000
13 4.0000 0.0000
14 -0.0000 5.0000
988 -0.0000 -5.0000
989 4.0000 -0.0000
998 0.0000 2.0000
So what is going on? On line 3 we define the sampling time dx in radians and also the time
x is specified in radians. Notice that we stop prior to 2 at 2 dx because of the periodic
assumption of the Fourier transform. On line 4 we define the test function, and on line 5 we
carry out the FFT. The output is in vector g and when you would inspect it you would see that
it contains complex numbers to store the Euler coefficients after the transformation. Also, the
numbering in the vector in MATLAB does matter in this discussion. At line 6 the indices in
the g vector are retrieved where the amplitude of the spectral line (defined as (A2i + Bi2 )1/2 )
exceeds a threshold. The FFT function is not per se exact, the relative error of the Euler terms
is typically around 15 significant digits which is because of the finite bit length of a variable in
MATLAB. If you find an error typically greater than approximately 10 significant digits then
inspect whether x is correctly defined. Remember that we are dealing with a periodic function f
and that the first entry in f (in MATLAB this is at location f(1)) repeats at 2. The last entry
in the f vector should therefore not be equal to the first value. This mistake is often made, and

43
it leads to errors that are significantly larger than the earlier mentioned 10 significant digits.
On line 7 the number of significant Euler pairs in the g vector are recovered, and on line 8 we
compute a scaling factor which is essential to understand what is stored in the g vector. The
part that decodes the g vector starts on line 9, integer i runs from 1 to n (the number of unique
pairs in g) and the scale factor is, depending on where we are in the g vector, adjusted on lines
10 to 13. The Euler terms for each spectral line are then recovered on lines 14 and 15 and the
result is printed on line 14. Line 15 terminates the for loop.
We learn from this program that vectors in MATLAB start at index 1, and not at zero as
they do in other programming languages. The value at g(1) = k A0 where k = N1 with N
denoting the number of samples on the input record f (and the definition of time in x). At the
4th index in g we find the complex number 0 2j = (0, 2) where j = 1, the sine term at the
third harmonic is therefore stored as (0, 2kB), at location 13 in g we see that the cosine term
is properly put at the 12th harmonic, it is stored as (2kA,0), location 14 in g confirms again
that the sine term at the 13th harmonic is stored as (0, 2kB). Next the g vector is completely
empty until we reach the end where we find the Euler coefficients stored in a reversed order
where the last term g(N ) contains k(A1 , B1 ), it is preceded by g(N 1) = k(A2 , B2 ) and so on.
To summarize the content of g after we executed g = fft(f ) in MATLAB:
1 2
First define a scaling term k = N for the zero frequency and k = N for all other frequencies.

The first location in the g vector contains the bias term: g(1) = k(A0 , 0)

g(i) for i > 1 and i < N/2 contains g(i) = k(Ai1 , Bi1 )

g(N i + 1) for i 1 and i < N/2 contains g(N i + 1) = k(Ai , Bi )

For this reason we say that the g vector is mirror symmetric about index N/2, and that the first
part of thevector contains the complex conjugate of the Euler coefficient pair A + jB = (A, B)
where j = 1. Furthermore the scaling term k should be applied. It also leaves one to wonder
what is going on at index N/2. In factor the sine term at that frequency evaluates as zero by
definition, so it does not exist.

4.3.2 Harmonics of a sawtooth function

The sawtooth function in figure 4.4 has a Fourier transform, and the question is asked, how
many harmonics to you need to approximate the function to 95% of its total power. You can do
this analytically with the help of the earlier integral definitions, but it is relatively easy to do
in MATLAB which is what we discuss hereafter. The function is shown in figure 4.4. In order
to solve this problem you need to do two things, first, compute the FFT of the input function,
next, check with the help of the Parseval theorem how much power is contained in the spectrum.
From the 0 (or DC or bias) frequency upward we will continue to look for the point where the
power contained in the lower part of the spectrum exceeds the 95% threshold which was asked
in the assignment. The result that I found is in figure 4.5. The conclusion is therefore that you
need at least 12 harmonics to reach 95% of the power contained in the input function. Lets go
over the MATLAB source to see how it is computed.

T = 1000; N=100; dt=T/N;

t = 0:dt:(T-dt);

44
Figure 4.4: The sawtooth function, also called the sweep generator function. Horizontal index
is time, vertical signal is output.

Figure 4.5: Remaining power contained of the sweep generator at a harmonic, it is expressed as
a percentage.

45
 x = zeros(size(t));
x = mod(4*t/T,1); k = 20;
figure(1); plot(t,x)
sum1 = sum(x.^2)/N; % sum in the time domain
X = fft(x)/N;
sum2 = abs(X(1)).^2 + 2*sum(abs(X(2:N/2)).^2); % sum in the spectrum
fprintf(Sum in the time domain is %15.10e\n,sum1);
fprintf(Sum in the freq domain is %15.10e\n,sum2);
fprintf(Relative error is %15.10e\n,(sum1-sum2)/sum1);
sum(1) = abs(X(1)).^2;
for i=2:N/2,
sum(i) = sum(i-1) + 2*abs(X(i)).^2;
end
percentage = (sum2-sum)/sum2*100;
harmonics = 0:N/2-1;
figure(2);
plot(harmonics(1:k),percentage(1:k),o-);
xlabel(Harmonics); ylabel(percentage power);
grid
After execution the program prints the message:
Sum in the time domain is 3.1360000000e-01
Sum in the freq domain is 3.1360000000e-01
Relative error is 0.0000000000e+00
The main steps in the program are that the function is defined and plotted on lines 1 to 5.
The power in the time domain is calculated in variable sum1, and the power in the spectrum
is collected in sum2, the following three print statements perfectly verify Parsevals theorem,
indeed, the power in the time domain is the power in the spectrum. No free energy here, why
should it exist anyway? After this step we compute the sum of the power in the spectrum
for each spectral line, this is the summing loop at lines 12 to 15, the percentage of what is
contained in the lower part of the spectrum relative to the total is then evaluated (it represents
a truncation error), next the results are plotted and the user is asked to find the point in the
graph where we go below the 5% point. This coincides at the 12th harmonics approximately.

4.3.3 Gibbs effect

The previous example is rapidly extended to demonstrate the so-called Gibbs effect (named after
its rediscoverer J. Willard Gibbs 1839 1903) which is a direct consequence of truncating the
spectral range of an input function. We could take for instance the function that we examined in
section 4.3.2 and examine the result after we truncate at the nth harmonics. More elegant is to
do this for the square wave function as is shown in figure 4.6. Obviously the resulting function
after band-pass filtering is distorted, the lower graph shows the typical Gibbs ringing at the
point where there is a sharp transition in the input function. It is relatively easy to explain
why we get to see a Gibbs effect after a Fourier transformation. The reason is a discrete input
signal sampled at N steps between [0, 2] can be represented as the sum of a number of pulse
functions that each come with a width t = 2/N . However, due to Nyquist we will also see

46
Figure 4.6: Top: square wave function, bottom: band-pass filtered version of the input function.
The Gibbs effect is now visible in the band-pass filtered result where you observe a ringing effect
at the points where there is a rapid (sharp) transition of the input signal.

that there is a maximum frequency. The only thing that we need to do is to position these pulse
functions at the right position (phase) along the time axis and to assign them an height as large
as the values that you find in the vn vector. Inevitably this will result in the computation of
the Fourier transform of a rectangular function, of which we know that the Fourier transform
1
will be sinc(x) = x sin(x). The sinc function will result in the ringing effect that we observe
at sharp transitions on the input signal.
In physics band-pass filtering is often the result of a physical property of a system. This is
relatively easy to demonstrate in the laboratory, take an oscilloscope to display the signal of a
square wave generator. If the frequency of the generated signal is well below the bandwidth of
the oscilloscope then sharp transitions are observed on the screen just like in the top graph of
figure 4.6. However, if we increase the frequency of the input signal then the oscilloscope will
display at some point the Gibbs effect, reason is that the bandwidth of the scope is insufficient
to deal with the highest harmonics that exist of the input function.

4.4 Power density spectra

Power spectra as discussed before should in fact be called power density spectra because the
frequency resolution is defined as f = 1/T with T being the record length. Each line i in the
spectrum should be associated with Pi = (A2i +Bi2 ), this part represents a part of the total power
over a limited slice f in the spectrum. Spectral density is too important to not represent it
along the vertical axis, because the natural question would then be to explain at the same time

47
what record length has been used in the frequency analysis. For this reason it is advisable to
represent the result as a power density, or the square root of a power density, because it would
be unambiguous to recover the power in the time domain without being dependent on the length
T of the data record during the analysis.
In a power density spectrum we therefore represent Pi /f along the vertical axis which has
the units [V ]2 /[Hz] if the input signal would be a voltage, thus in units of [V] which was sampled
over a certain length in time. An integral over the frequency in the power density spectrum
would in the end recover the power in the time domain, this could be the total power in the
time domain, or it could be the power of a band-pass limited version of the signal in case we
decide to truncate it. Sometimes the square root of the power is displaced along the vertical
axis while it is still a density spectrum. In the latter case we find the units [V ]/[ Hz] along the
vertical axis in the spectrum. Sometimes alternative representations than the Hertz are used
and spectra are represented by for instance wave-numbers.

4.5 Exercises
Here are some examples:

1. Apply a phase shift in the time domain of the test function in eq. (4.18) and verify the
results after FFT in MATLAB. To do this properly you compute the function f (x + phi)
for a non-trivial value of phi in radians. In the time domain this results in a new function
definition where you are able to compute the amplitudes and phases at each spectral line,
the same result should appear after FFT. This test is called a phase stability test, is it
true, or is it not true?

2. Implement the convolution of the f and the g block functions as shown in section 4.2.3 to
recover the h function in MATLAB. What are the correct scaling factors to reproduce h?

3. Verify that the Euler terms of a square wave function match the analytical Euler terms.
In this case you can use MAPLE to derive the analytical Euler terms, and MATLAB to
verify the result.

4. Take the solar flux data from the space weather website at NOAA (or any other source).
Select a sufficient number of years and find daily data. Where is most of the energy in the
spectrum concentrated. Apply a tapering function to the data and explain the difference
with the first spectrum.

5. Demonstrate that you get a Gibbs effect when you take the FFT or a sawtooth function,
how many harmonics do you need to suppress the Gibbs effect?

48
Chapter 5

Reference Systems

Within any laboratory the definition of coordinates and time is artificial, we can chose an
arbitrary coordinate offset and orientation or set the reference time to an arbitrary epoch, any
choice of these settings should be inconsequential for the measurements within our laboratory,
and the choice of the reference should not affect the experiment that we carry out within the
framework of our laboratory. However, if an experiment within the walls of the laboratory
depends on the outcome of an experiment in another lab, then the transfer of information from
one reference system to another becomes relevant. The meter could for instance be defined in one
lab, and it should be the same meter that we use in the other lab in order to be consistent. As
soon as this discussion will affect the outcome of an experiment then reference system information
needs to be exchanged in a consistent way. For this reason there is the International System
of Units, (French: Systeme Internationale: SI) that provides the best possible definitions of
reference system information relevant to accurately execute an experiment in your laboratory.
When we work with coordinates and time related to satellites technique the same discussion takes
place. The scale of the laboratory is much larger, in principle it extends to the dimensions of
planet Earth (or the solar system) for which a reference for coordinates and time must be chosen.
This chapter is about the choice that reference, but also the choice to convert experiments from
one reference system into another reference system.
Where do we start? The Kepler problem discussed in chapter 2 gives us the opportunity to
set-up a reference system, at the same time, potential functions that describe the gravity field
of a planet ask for a definition of a reference system. Whereas the definition of time and
coordinates is perhaps the most fundamental definition of a reference system, some people also
take the definition of models as a subject for a reference system. In the following we will start
with a few examples and extend the definition to the space geodesy where it is necessary to
relate coordinates and time on the Earths surface to coordinates and time for satellites orbiting
the Earth, or any other planet or moon in our solar system. The discussion of a coordinate base
not only finds its application in two dimensions on a national or international level, in fact, in
any construction workshop you can do the same, so that coordinates in the lab are uniquely
defined.
The simplest way to define a reference system in two dimensions is to assume a plane with
at least one reference point, by definition one coordinate consists of two ordinates. Next we
need an orientation of the x-axis and a measure of length, a so-called scale in the reference
system. But there are other possibilities, we could define two reference points in the plane so
that orientation, scale and origin are defined. The degrees of freedom in this problem is 4 when

49
you work with planar coordinates. In a three dimensional space the definition of coordinates is
less obvious. Possible solutions are two reference points and one extra direction to define the
third axis. But another possibility is one origin, two directions and one measure of length, and
a third direction to complement the frame. No matter what we do, a the three dimensional
reference system has seven degrees of freedom and those 7 numbers need to be defined.
Intermezzo: Within the Netherlands, as well as many other countries, surveying networks
can be connected to a coordinate base. Before GPS was accepted as a measurement technique
for surveying there was a calibration base on the Zilvensche heide in the Netherlands. For more
information see https://nl.wikipedia.org/wiki/IJkbasis.
The next problem is that we are dealing with two applications for coordinates, namely,
coordinates of objects attached to the surface of a planet or moon in the solar system, or,
coordinates that should be used for the calculation of satellite trajectories where we want that
Newtons laws may be applied. Within the scope of orbit determination it is not that obvious
how we should define an inertial coordinate system. We may either chose it in the origin of
the Sun, or the Earth, or maybe even any other body in the solar system, but for Earth bound
satellites we speak about an Earth Center Inertial (ECI) system. Within the scope of tracking
systems on the Earths surface we assign coordinates that are body fixed, this type of definition
is called an Earth Center Fixed (ECF) system. The relation between the ECI and the ECF
system will be discussed in section 5.1 and the definition is further worked out in section 5.1.1.
Another issue is that ECF coordinates may be represented in different forms, we could choose
to represent the coordinates in a cartesian coordinate frame, or, alternatively, we may choose
to represent the coordinates in a geocentric or a geodetic frame. Furthermore coordinates are
often represented as either local coordinates where they are valid relative to a reference point
or they may be represented globally. The ECF coordinate representation problem is discussed
in section 5.2.
The definition of time should also be discussed because, first there is the problem of the
definition of atomic time systems in relation to Earth rotated and the definition of UTC, this
is mentioned in the context of the IERS, see section 5.1.2, which is the organization responsible
for monitoring Earth rotation relative to the international atomic time TAI. For the definition
of time also relativity plays a role, and this topic is discussed in section 5.3.

5.1 Definitions of ECI and ECF

For orbit determination within the context of space geodesy involving satellites near the Earth
specific agreements have been made on how the ECI system is defined. Input for these definitions
are the Earths orbital plane about the Sun, the so-called ecliptic, and the rotation axis of the
Earth, and in particular the equatorial plane perpendicular to the Earths rotation axis. For
the ECI frame the positive x-axis is pointing towards the so-called vernal equinox, which is the
intersection of the Earths equator and the the Earths ecliptic. The z-axis of the Earths inertial
frame then points along the rotation axis of the Earth. In [63] this explained in section 2.4. This
version is called the conventional inertial reference system, short: CIS in some literature, or the
Earth centered inertial frame, the ECI in [63]. All equations of motion for terrestrial precision
orbit determination may be formulated in this frame. The ECI should be free of pseudo forces,

50
so that the equations of motion can assume Newtonian mechanics. 1
For the ECI we defined also 7 parameters. The first assumption is that the ECI frame is
centered in the Earths origin (3 ordinates), the direction toward the astronomic vernal equinox
and the orientation of the z-axis defined (in total 3 directions), and the scale of the reference
system is the meter. For the ECF system the situation is similar, in this case the coordinates
are body-fixed, and several rotations angles are used to connect the ECI to the ECF.

5.1.1 Transformations
The transformation between the ECI and the ECF is:

xECF = SN P xECI (5.1)

where S, N and P are sequences of rotation matrices.

S = R2 (xp )R1 (yp )R3 (GAST )

N = R1 ( )R3 ()R1 ()
P = R3 (z)R2 ()R3 ()

and where GAST = GM ST cos  describes the difference between the Greenwich Ac-
tual Sidereal time (GAST) and the Greenwich Mean Sidereal Time (GMST). The difference is
described by the so-called equation of equinoxes which in turn depends on terms that one
encounters within the nutation matrix. We remark that:
The precession effect is caused by the torque of the gravity field of the Sun on an oblate
rotating ellipsoid which is to first order a good assumption of the Earths shape. The
Earth rotation axis is perpendicular to the equator, and the equatorial plane is inclined
with respect to the ecliptic. The consequence is that the Earths rotation axis will precess
along a virtual cone, a characteristic period for this motion is approximately 25772 years.
To calculate the precession matrix P we need three polynomial expressions, details can be
found for instance in [60] eq.(2.18). One should be careful which version of the precession
equations are used because different versions exist for the ECI defined at epoch 1950 and
the ECI at epoch 2000. In literature these systems are called J1950 and J2000 respectively.
Furthermore the precession effect of the Earth hardly changes within a year, therefore the
choice is made in numerous software packages to calculate the P matrix only once, for
instance in the middle of a calculated satellite trajectory.

Another effect that is part of the transformation between the systems concerns the nuta-
tion effect, which is in principle the same as the precession effect, except that the Moon
is responsible for the torque on the Earths ellipsoid. The N matrix is far more costly
to compute because the involved nutation angles consist of lengthy series expansions with
goniometric functions (sin and cos functions). Within most programming languages go-
niometric functions are evaluated as polynomial approximations, that these calculations
are by definition expensive. Also in this case it is desirable to compute the N matrix once,
and to leave an approximation of the N matrix in the calculations.
1
A pseudo force is perhaps a bit of a strange concept, you might have experienced it as a child sitting in the
center of a spinning disc in the playground. Sitting there in the center way fine, but dont try to go from the
center to the edge because the Coriolis effect will cause you to fall.

51
 Within the S matrix we encounter the definition of the GMST which says in essence
that the Earth rotates in approximately 23 hoursand 56 minutes about the z-axis of the
ECF frame. The equation for the GMST angle follows for instance from equation (2.49)
in [60], it is a compact equation and it is relatively cheap to evaluate it quickly. The
difference between GMST and GAST is a slowly changing effect whereby the definition of
the nutation matrix is relevant. The GMST variable must be computed in the UT1 time,
and not the leap second corrected UTC time system that we may be used to for civil
applications. The International Earth Rotation Service, the IERS, is the organization
responsible for distributing the leap second, more on this part will follow later in this
chapter. The remaining effects in the S matrix are the polar motion terms xp and yp , also
these terms are disseminated by the IERS. The values of xp and yp are in units of milliarc
seconds, and they follow from analysis of space geodetic observations.

5.1.2 Role of the IERS

As was explained before, for the S matrix we need three variables, xp and yp and the difference
between UT1 and UTC (short UT1-UTC) because the observed or computed time (specified
in UTC) needs to be converted to UT1 known as Earth rotation time. The variables in S
are available for trajectories before the present, but, there is no accurate method to predict xp ,
yp and UT1-UTC a number of weeks ahead in time. The International Earth Rotation Service
(IERS) is established to provide users access to xp , yp and U T 1 U T C. They collect various
estimations of this data and have the task to introduce leap seconds when |U T 1 U T C| exceeds
one second. The IERS data comes from various institutions that are concerned with precision
orbit determination (POD) and VLBI, and collects summaries of the different organizations
including predictions roughly a month or so ahead in time of all data. For precision POD
predictions are not sufficient, and use should be made of the summaries for instance in the IERS
bulletin Bs. For all precision POD applications this means that there is a delay usually as
large as the reconstruction interval of one month that the IERS needs to produce bulletin B.
The predicted IERS values are of use for operational space flight application, for instance in
determining parameters required for launching a spacecraft to dock with the international space
station. In the past bulletin Bs were sent around by regular surface mail, currently you retrieve
them via the internet.

5.1.3 True of date systemen

In literature we find the terminology true of date (TOD) to specify a temporary reference
system. TOD systems are used to make a choice for a quasi inertial reference system that differs
from J2000. For realizing a TOD system the P and N matrix in (5.1) are set to a unit matrix,
precession and nutation effects are not referring to the reference epoch of J2000, but instead a
reference time is chosen that corresponds to the current date, hence the name True of Date.
All calculations between inertial and Earth center fixed relative to such a chosen date should
not differ too in time much relative to this reference date. The benefit of TOD calculations
is that the P and the N matrix dont need to be calculated at all epochs, so this saves time.
However, the S matrix does need frequent updates because the involved variables, GAST , xp
and yp change more quickly. For POD to terrestrial satellites whereby the orbital arc does not
span more than a few days to weeks the accuracy of the calculations is not significantly affected

52
by assuming a TOD reference system where afterwards satellite state vectors in the TOD system
are converted to J2000.

5.2 Representation problem

Coordinates in an ECF frame may be defined in the way the IERS recommends to implement
an international terrestrial reference frame (ITRF) for which different versions exist. Essentially
the ITRF is maintained by providing a list of cartesian coordinates for a number of tracking sta-
tions. Cartesian refers to the French mathematician Rene Descartes 15961650 who introduced
Cartesianism in mathematics. Coordinates of tracking stations in a reference system may be
represented as cartesian coordinates, but the reality is that also other representation forms are
used to denote the same points. Alternative ways to represent coordinates are discussed in the
following subsections.

5.2.1 Geocentric coordinates

The relation between cartesian coordinates (x, y, z) and geocentric coordinates (r, , ) is:

x = r cos cos (5.2)

y = r sin cos (5.3)
z = r sin (5.4)

where (r, , ) denote the geocentric coordinates radius, latitude and longitude. If (r, , ) are
known then the equations implement the transformation, the inverse transformation is also
known although singularities exist at both poles. Geocentric coordinates are used for instance
to obtain the spherical coordinates required for a Legendre function expression to obtain the
gravitational potential and its derived quantities.
The local variant of geocentric coordinates also exists, in this case the local coordinate frame
is taken at the location of the observer, the x-axis is pointing to the local East, the y-axis is
pointing to the local north, and the z-axis is then taken along the vertical. In case we plot the
vector from the observer to an object in the sky we speak about topocentric coordinates. The
azimuth is the angle relative to the north, and the elevation is the angle relative to the local
horizon. The lecture slides contain images of geocentric and topocentric coordinates.

5.2.2 Geodetic coordinates

The relation between Cartesian coordinates (x, y, z) and geodetic coordinates (r, 0 , ) may only
be understood if we introduce a reference ellipsoid, see figure 5.1. A reference ellipsoid is a
mathematical figure that defines the best fitting mathematical figure of the Earths sea level at
rest so that the sum of all geoid heights relative to the reference ellipsoid is minimized. Not
only is the reference ellipsoid a best fitting surface, it is also a body on which the gravity field
should be defined. Note that there is a difference between gravity and gravitation, whenever
gravitation is mentioned Newtons gravity law should be put in place, but when you talk about
gravity then also the centrifugal effect of a rotating body should be taken considered. In the
end there are four parameters that define the reference ellipsoid:

53
Figure 5.1: Relation between geocentric and geodetic coordinates on a reference ellipsoid.

We need the shape of the body that is represented in the parameters Re and Rp which
are the equatorial and the polar radius, you will find them in eqns. (5.5) where we use the
parameters Re and f , the latter is a derivative of Re and Rp .

The total mass of the body, this is represented in the parameter which is the gravitational
constant G measured in the laboratory, times the mass of the Earth M . The product
follows from orbit determination, in fact, the value of is better determined than G or
M individually. Error propagation shows that a n
= 3 a + 2 n , the relative error in

is determined by our ability to determine the orbital period and the semi-major axis of
a satellite. After sixty years of orbit determination we presently know to be equal to
3.986004418 1014 (4 105 ) m3 s2 , so the relative error is of the order 1 part per billion
(ppb), whereas the relative error in G = 6.67408(31) 1011 m3 kg 1 s2 is of the order of
10 part per million, hence 10000 times less accurate, essentially because the metrology of
the experiment in the laboratory determines the number of significant digits by which we
know G. The value of M for the Earth is inferred from and G, so its relative error is
never better than GG .

When only of the reference ellipsoid is defined we can not yet define gravity on its surface,
for this you also need the rate of rotation of the ellipsoidal figure in inertial space. This
parameter could be called and it should not be confused with the argument of perigee
of a satellite orbiting the Earth. The value of depends on our ability to measure the
length of a sidereal day, since it is a time measurement it is relatively easy to do, relative
accuracies of 1 ppb are easily achieved, typically we get 1 micro-second in error over a

54
 Figure 5.2: Geoid relative to reference ellipsoid

length of a day,
1011 as far as the measurement accuracy is concerned.

With the above definition of a reference ellipsoid we arrive at the transformation between geo-
centric and ellipsoidal coordinates, the relation is:

x = (N 0 + h) cos cos 0
y = (N 0 + h) sin cos 0 (5.5)
0 2 0 0
z = (N + h e N ) sin

where (0 , , h) denote the geodetic coordinates latitude, longitude and height above the reference
ellipsoid. The geodetic height above the reference ellipsoid depends on a number of ellipsoidal
parameters, namely the semi-major axis Re and the semi-minor axis Rp which are used to
R R
defined the flattening parameter f = eRe p . Parameter N 0 = (1e2 sin Re
2 0 )1/2 depends in turn

on e = f (2 f ). For ellipsoidal coordinates on the reference ellipsoid we know that x2 + y 2 +

(R p
) z = Re2 . The conversion of (h, 0 , ) to (x, y, z) is straightforward, the inverse relation has
Re 2 2

no easy analytical solutions, iterative methods are used for this purpose.

Geoid height
In the previous section we mentioned the concept geoid heights which we call N in the sequel.
The geoid in figure 5.2 is defined as the imaginary surface that coincides with the mean sea
level at rest, hence, there are no currents or tides and the shape is that of an equipotential
surface which follows from the definition of the geo-potential V which we derived in eqn.(3.3).
Topographic heights (land, hill, mountains) are normally presented relative to the geoid, and the

55
geoid in turn is defined relative to the reference ellipsoid. Bathymetric heights (sea floor, ocean
depth), could be represented to the geoid, but, hydrographic services chart the bathymetry
relative to the lowest possible tidal surface because their interest is in safe traffic over coastal
seas.
In order to obtain N in figure 5.2 you need a model for the Earths gravity field, and also
you need a model that defines the reference ellipsoid. In orbit determination there are analysis
groups that determine series of so-called potential coefficients C nm and S nm that follow from
the Cnma coefficients in eq. (3.16). The geo-potential field V of the Earth is then defined as:
N n 
X X Re n+1 


V (r, , ) = + C nm cos(m) + S nm sin(m) P nm (sin ) (5.6)
r Re r
n=2 m=0

This equation is established in a geocentric Earth fixed reference frame (there are expressions for
ellipsoidal harmonics, but we dont need them for this lecture) that comes with the geocentric
coordinates r, and . The maximum degree n and order m expansion of the geo-potential
V of the Earth is currently known up to degree and order 240 since that the GOCE gravity
gradiometer mission mapped the field. From satellite orbit determination methods we presently
know the geo-potential V to N 90. The GRACE satellite to satellite tracking mission produces
monthly maps of V up the N = 60.
The higher the satellite altitude above the Earths surface, the fewer details of the geoid (or
gravity anomalies) are sensed by precision orbit determination methods, with the Lageos mission
orbiting the Earth at around 5900 km above its surface we can determine gravity field details to

n+1
N 6. The main reason is that the term Rre results in a damping reducing the sensitivity
of the orbit determination method to determine the higher degree and orders. With the help
of satellite altimetry we can directly measure the gravity anomalies g at the sea-surface, the
resolution is in this case phenomenal, typically better than 10 km, but the spatial coverage is
confined to the oceanic domain.
In order to compute a geoid height N on the surface of the Earth we consider the ellipsoidal
model that we introduced in section 5.2.2. The reference ellipsoid parameters allow one to
compute a constant value of the geo-potential along the reference ellipsoid. What the reference
ellipsoid parameters also do is that they allow one to compute a normal field expression for the
gravity acceleration at the reference ellipsoid. If the normal field expression of the reference
ellipsoid is called W , then N follows from the Bruns equation:
T
T =V W N= (5.7)

where T is called the disturbance potential. The relation between the earlier mentioned gravity
anomalies g and the disturbance potential T is:
T 2
g = T (5.8)
r r
The term T r is called the gravity disturbance g. Physical geodesy is the study of the potential
field of the Earth, the so-called geodetic boundary value problem may follow from equation (5.8),
thus g is given along the physical surface of the Earth, thus at the topographic height H relative
to the geoid N in a defined reference system, and the task is to solve for the disturbance potential
T . Various solutions for this problem were proposed prior to the era of spaceflight, one of the

56
possible solutions is the solve for the geoid height N via a so-called Stokes integral over a field
of observed gravity anomalies: Z
Re
N= g St () d (5.9)
4
where it is assumed that the gravity anomalies are observed on the geoid, and where St () is
the so-called Stokes function, for details see [26].
The same technique of gravity field determination and reference ellipsoid estimation can be
established on other planets and moons in the solar system. The MOLA laser altimeter from
NASA that orbited Mars has resulted in detailed topographic maps and representations of the
geoid. From such information we can learn a lot about the history of a planetary surface, and
the internal structure of the planet. On Earth we confirmed the existence of plate tectonics by
satellite methods, the gravity feature of plate tectonics was earlier discovered by Felix Vening
Meinesz who sailed as a scientific passenger with his gravimeter instrument on a Navy submarine.
Currently we know that planet Earth is probably the only planet where plate tectonics exist,
Mars does not show the features of plate tectonics in its gravity field although magnetometer
mapping results do seem to confirm some tiger stripes typical for plate tectonics. Venus would be
another candidate for plate tectonics, it was extensively mapped by NASAs Magellan mission
but also here there is no evidence for plate tectonics as we find it on Earth.

5.2.3 Map coordinates

Coordinates on the surface of the ellipsoid may be provided on a map which is a cartesian
approximation of a part or the entire domain. This is a cartographic subject that we do not
work out in these notes, instead the reader is referred to [61]. Well known projections are the
Mercator projection, Lambert conical, UTM and the stereographic projection. There are also
more eccentric projections like that of Mollweide which simply look better than the Mercator
projection where the polar areas are magnified. Topographic coordinates have a military appli-
cation, because azimuths found in the map are equal to the azimuth found in the terrain which
aids navigation and targeting.

5.3 What physics should we use?

Is Newtonian physics sufficient for what we do, or, should the problem be extended to general
relativity? For special relativity the question seems to be relevant because we are dealing with
velocities between 103 to 104 meters per second relative to an observer on Earth. Furthermore
Earth itself has a rotational speed of the order of 2.87 104 m/s relative to the Sun, and the
Sun has a rotational speed relative to our galaxy.
For special relativity the square of the ratio of velocity to the speed of light becomes relevant,
thus (v/c)2 so that the scaling factors become approximately 108 for time and length. For
general relativity another effect becomes relevant, in this case the curvature of space and time
caused by the gravity field of anything in the solar system needs to be considered. All masses
generate a curvature in space and time, for our applications Earth and Sun seem to be the most
relevant masses. Time-space curvature turns out to be relevant in the definition of reference
systems and in particular the clock corrections that we will face in the processing of the data.

57
 In the case of radio astronomy, and in particular VLBI, the change in the direction of
propagation of electromagnetic waves is observable near the Sun2 . There is quite some literature
on the topic of general relativity, the reader may want to consult [62] but also [48]. Within the
scope of these lecture notes I want to discuss time dilatation and orbital effects that affect the
clocks and orbits. Also I want to spend some time on the consequence of general relativity on
clocks.

5.3.1 Relativistic time dilatation

Time is presently monitored by a network of atomic frequency standards that have a frequency
stability far better than one second in a million year equivalent to (f /f ) < 3 1013 where f
is the frequency of the clocks oscillator. To understand relativistic time dilatation one should
distinguish between two observers, one on the ground and one on a satellite. For the terrestrial
observer it will appear (within the framework of special relativity) as if the satellite clock is
running slower compared to his clock on Earth. Why is this the case? Albert Einstein who
came up with these ideas introduced the assumption that the speed of light c is independent of
the choice of any reference system. So it would not matter for a moving observer to measure
c in his frame, or for an observer on Earth to measure c, in both cases they would get the
same answer. The assumption made by Einstein was not a wild guess, in fact, it was the most
reasonable explanation for the Michelson-Morley experiment whereby an interferometer is used
to detect whether Earth rotation had an effect on c. The conclusion of the experiment was that
it did not matter how you would orient the interferometer, there was no effect, see also chapter
15 in the Feynman lecture notes [48].

Intermezzo
Suppose that we align two mirror exactly parallel and that a ray of light bounces between both
mirrors. If the distance between both mirrors is d then the frequency of the light ray would
c
be f = 2d . So if d is equal to e.g. 1 meter then f = 150 MHz which is just above the FM
broadcast spectrum. Suppose now that we construct a clock where this light oscillator is used
as the reference frequency. Electronically we measure the frequency, and we divide it by 150
million to end up at a second pulse. This pulse is fed into a counter and this makes our clock.
The light-clock is demonstrated in figure 5.3, in the left figure the light travels between A and
B along the orange dashed line.
Now we add one extra complication, we are going to watch at the light clock where both
parallel mirrors move along with a certain speed v as is shown in figure 5.3 in the right part.
For an observer that is moving with the experiment there is no problem, he will see that the
light ray goes from one mirror to another, and back, thus like in the left part of figure 5.3. The
speed of the right ray will be c according to Einsteins theory of relativity. This was also found
with the Michelson-Morley experiment, so for an observer who travels with the reference frame
of the interferometer there is no effect of v on the outcome of c.
But lets now look from the point of view of an observer how watches the light clock from a
distance, thus outside the reference frame of the light clock. For the stationary observer it will
appear as if the light ray starts at A in figure 5.3 that it travels to B along the red dashed line,
2
In essence this is a variant of the proof of validity of the theory of general relativity where the perihelium
precession of the planet Mercury was observed.

58
Figure 5.3: Light-clock experiment, left: two parallel mirrors in the light clock where a light
beam is bouncing between the mirrors, right: the same experiment where the light-clock is seen
from a stationary observer

and that it returns to C. He will still see that the light ray travels at speed c, but, the length of
the path has increased so that it needs more time to travel up and down between the mirrors,
or a reduced speed as indicated by the green dashed
path between B and D as in figure 5.3. For
the external observer the reduced frequency is c2 v 2 and he will see a frequency f like:

c2 v 2 c p p
f = = 1 (v/c)2 = f 1 (v/c)2 (5.10)
2l 2d
We see that f is scaled with respect to f , and this effect is called time dilatation within the
context of special relativity. Moving clocks run slower for stationary observers is the main
summary of the theory. For the above example where d = 1 meter and v = 103 we find that
f f = 0.9 103 Hz. This effect seems small, but, is may be detected easily with modern
atomic time standards. The time dilatation is 480 nsec per day or -175 sec per year.

Lorentz contraction
Apart from time-dilatation there is also length-contraction or Lorentz-contraction if we assume
that the speed of light c does not depend on the velocity v of the reference frame. Also this
is easily understood, in particular when we take the velocity component v parallel to the line
DB in figure 5.3. For an observer who moves along with the mirrors the time it takes to moves
between D and B has not changed, it remains T = 2d c . But a stationary observer, who watches
the experiment from a certain distance, will notice a slower oscillation with a periodicity T , the

59
period will not be different from what we already found for time dilatation. The consequence is:
p
T 2l 2l 1 (v/c)2 2l
T =p = T = =
1 (v/c)2 c c c

from which we see that:

l = l
p
1 (v/c)2
The conclusion is that objects in rest will have a length l, but when they move relative to
an observer it will appear as if they become shorter. For completeness we show the complete
Lorentz transformation where both length and time are transformed:
x vt
x0 = p ,
1 (v/c)2
y0 = y,
z 0 = z, (5.11)
t vx/c2
t0 = p
1 (v/c)2

This transformation applies between the (x, y, z, t) system and (x0 , y 0 , z 0 , t0 ) system for the rela-
tively simple case where two observers have a relative motion with velocity v along a common
x direction, see also [48]

5.3.2 Gravitational redshift

Apart from time dilatation and Lorentz contraction within the context of special relativity there
is a relation between the position within a gravity field and the rate of a clock oscillator. This
problem is called the gravitational red-shift problem, which we put under the heading of the
general theory of relativity. Figure 5.4 shows a local reference frame near a star. A photon is sent
away from the star and it has a certain color that matches frequency f as indicated in figure 5.4.
The photon can only fly at the speed of light c, and, the gravity g of the star is now supposed
to affect the photon. How can it do that? If the photon had a mass, then you would expect
that it slows down in the presence of the gravity of the star, in that case the change of velocity
dv in a time interval dt would be dv = a.dt where a is the inertial acceleration experienced by
the particle. And if we assume that the equivalence principle3 is valid, then the acceleration
experienced by the particle would be equal to the gravitational acceleration (we called that the
gravity) of the star. If the particle had traveled over a distance dh then dv = g.dt, and therefore
the change in velocity is dv = g dhc .
A property of photons is that they can not change their velocity or their mass. Photons (in
vacuum) travel at the speed of light c without any mass. All energy in the photon goes into
its frequency f and for this there is Plancks equation E = h.f where h is Plancks constant.
To change the energy of the photon we can however change its frequency. The dv that we had
3
The equivalence principle follows from the tower experiment in Pisa, where one has seen that the acceleration
experienced by a mass does not depend on the mass of the ball thrown from the tower itself. Both balls did hit
the ground at the same time, and as a result inertial mass is equivalent to gravitational mass. In other words,
2
any mass term in f = m.a is equivalent to the mass term in Newtons gravity law where f = (Gm1 m2 )/r12

60
 Figure 5.4: The gravitational red-shift experiment

obtained before can be used for this purpose, we insert it in the Doppler equation and compute
a frequency shift df , so this is what happens:
 
dv
f + df = 1 + f
c
where we substitute dv = g.dh, so that
df dv g.dh
= = 2
f c c
which is a good approximation when variations in g are small, for larger values of dh we should
use the expression:
df d
= 2 (5.12)
f c
where the star is represented as a point mass so that its potential (r) = r where r is the
distance relative its center and its gravitational constant. The consequence is that clocks will
run at a different rate when you place them at different positions around the star. The more
heavy the star is, the more its emitted photons will be shifted to lower frequencies when the star
is seen at a distance, thus the Fraunhofer emission lines of the atoms responsible for generating
the starlight automatically move to the red part of the spectrum. For this reason we speak
about the gravitational redshift problem.
Related to the redshift problem is the Olbers paradox: Why is the nightly sky mostly dark
and why do we see only some stars? If the universe is infinitely large then starlight should
be seen in every corner of the skies. We dont see it so what is going on? The mainstream
explanation is that the gravitational redshift of all the starlight, and in particular the light

61
since the Big Bang, is shifted to the red. In the end cosmic background radiation with a
temperature of 2.76K remains. Maps of the CBR has been made with the COBE mission, see
for instance http://www.nasa.gov/topics/universe/features/cobe 20th.html where you
find how temperature differences in the CBR are measured by COBE.

Example
In a network of atomic frequency standards we have to account for the height of the clock
relative to the mean sea level, evidently, because f
f depends on the position of the clock in the
potential field, here, the altitude of the clock. Suppose that the network consists of a clock in
Boulder Colorado at 1640 meter above the mean sea level, while another clock in Greenwich UK
at 24 meter above the mean sea level. What is then the frequency correction and the clock drift
for the Colorado clock to make it compatible with the one at Greenwich? For this problem we
need the gravity potential as a function of the height h so that d can be computed in the
following way for both clocks:

b ghb = 9.81 1640 = 1.609 104 J/kg

g ghg = 9.81 24 = 2.354 102 J/kg

In other words:
d = (b g ) = 1.585 104 J/kg
From which we get:
df d
= 2 = 1.76 1013
f c
so that the clock in Boulder needs a negative correction to make it compatible with a clock in
Greenwich, the correction is 1.76 1013 86400 365.25 or 5.6 sec per year.

5.3.3 Schwarzschild en Lense-Thirring on satellite orbit dynamics

In [62] three relativistic are mentioned that play a role in the equations of motion of a satellite
in a gravity field of a planet. The gravitational redshift effect, or the Schwarzschild effect,
add terms to the equations of motion, so does the Lense Thirring effect and a smaller effect
mentioned in [62]. The consequence is that it appears as if large orbit perturbations will appear
in a satellite trajectory that you otherwise calculated with a Newtonian model.
The largest effect is caused by the Schwarzschild effect and the perturbations grow in par-
ticular in the direction of motion of the spacecraft. A deeper investigation on the cause of this
effect shows that it is not that relevant. It turns out that the same relativistic effect in the
orbit could also be obtained by scaling the gravitational constant = G.M of the planet in
the software program for orbit calculations that was originally based on the Newtonian model.
This conclusion should be seen as a general warning for relativistic orbit effects, a scaling of
some the used constants in a Newtonian reference model will usually result in the same effect
as applying a relativistic term in the equations of motion. The only reason to implement the
relativistic orbit perturbation is then to be consistent with the organization that provided you
with reference constants for your calculations.
The Lense-Thirring effect described in [62] leads to a different type of orbit perturbation.
The consequence of the LT effect is that it changes the precession of the orbital plane. The

62
same effect may be obtained for a single satellite via the gravitational flattening term J2 . A
satellite test of the theory of relativity has been performed with the Lageos experiment whereby
Lageos 1 and 2 were put in complimentary orbital planes. In this way the Newtonian effect of
gravitational J2 precession on the satellite orbit becomes irrelevant so that the LT effect becomes
visible; the dual plane approach was successfully applied for this purpose, see for instance [9].
This experiment ran in close competition with the GP-B mission designed by the university of
Stanford. Its purpose was also to test the same hypothesis of the theory of general relativity
by direct measured of the Lense-Thirring effect. GP-B contained a cryostat, a telescope and a
gyroscope and it was active between April 2004 and August 2005, the results of the mission are
described in [20].

5.4 Exercises
1. We determine the shape of a network of stations in 2D with the help of distance and
angle measurements. Generate a network of 5 stations and connect them. Show that the
measurements are unaffected when we rotate, or translate the coordinates.

2. Take three points on the sphere, distance measurements are now along great circles, and
angular measurements are as in spherical trigonometry. How do you compute the distances
and the angles when the coordinates of the points are known? Also invert the problem,
if you had the measurements, then what could you do with the coordinates to leave the
measurements as they are.

3. Which variables define the transformation between Earth Center Fixed and inertial coor-
dinates. Not only just mention the variables, but also explain the physics of what they
represent

4. Within the International Space Station a researcher creates his own time system which
is based on an atomic clock. There is a similar atomic clock on the ground and the
information of both clocks is required for an experiment. Explain how the researcher in
the space station needs to apply relativistic corrections to the clock on Earth to use its
information into his reference system.

5. Research has shown that atmospheric winds result in an angular momentum that interacts
with Earth rotation. Explain which variables are likely to be affected in the transformation
between terrestrial and inertial coordinates.

6. Describe the parameters that define the reference ellipsoid on Earth. What type of satellite
measurements are necessary to determine similar reference parameters on another planet
or moon in the solar system?

7. How do we compute a geoid height on Earth? Explain all assumptions that are made, also
assume that satellite tracking methods provided you the potential coefficients.

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Chapter 6

Observation and Application

Satellite laser ranging, doppler tracking and the global positioning system are nowadays in one
way or another used for the estimation of parameters when we determine a satellite trajectories.
All techniques described in section 6.1 obtain range measurements with or without a bias, or
the change of a distance along the line of sight between two epochs. A few typical corrections
are described in section 6.2, either these corrections deal with physical properties of the Earths
atmosphere, the ionosphere, or the finite speed of light, see also section 6.2. Scientific applica-
tions are discussed in section 6.3, here we mention very long baseline interferometry, satellite
altimetry and satellite gravimetry.

6.1 Techniques
We send a light ray from a transmitter to a receiver and measure the time difference. The
distance between both follows from the speed of light multiplied by the time difference between
both clocks. In figure 6.1 this configuration is indicated as A, it is the configuration that you
could chose in the laboratory. This is a somewhat straightforward experiment since we can use
the same clock, synchronization is realized by cables in the laboratory. When transmitter and
receiver are separated by a greater length this would become impractical, for these configurations
we go to set-ups B and C in the figure.
In configuration B we put the transmitter, the receiver and the clock in the same instrument,
and we added a reflector. This is a more practical configuration because it puts all the complexity
on one side. Reflectors can be constructed for both radio and laser techniques, they take the
shape as drawn in the figure, think of a ball returning to where it came from on a pool table
hitting a right angled corner. In three dimensions the same properties are found in the corner of
a glass cube. In daily life cats eye reflectors are found along roads, airport runways, the shore
and river banks where they facilitate navigation.
We also introduce a third maybe less obvious configuration, set-up C in figure 6.1 where
we use non-synchronized clocks both at the transmitter and the receiver. In the end the user
gets access to all information, thus when the signal left and it arrived. We can easily add
more transmitters and receivers to this configuration, it does not matter for the principle. An
important property of experiment C is that the range measurement will be biased as a result
of differences between the involved clocks. These biases will also exist for experiment A and
B; in the sequel we assume that a calibration techniques exists to eliminate these measurement

64
Figure 6.1: Three conceptual systems to measure distances between transmitters and receivers.
In all cases there are clocks to measure the transmit and receive time. The green line is an
idealized path of an electromagnetic wave, it can be laser light or a radio signal.

65
biases. The clock problem in set-up C seems like a significant drawback, but there are also ways
to get around it, these techniques are discussed in the section on GPS.
What also matters for all configurations are deviations from the theoretical speed of light in
vacuum due to refraction. These differences can be ignored as long as distances are not too long,
up to a few hundred meter ranges can be measured to within a few mm, the technique is used
for instance in surveying, but also in civil engineering, geophysical prospecting, construction
engineering and aviation navigation systems. Let us now scale the experiment and implement
measurements between ground and a satellite, or the other way around, between satellite and
ground, it does not really matter for the experiment. We assume in all cases that there is
approximate knowledge of ground station coordinates and the satellite trajectory, and that the
purpose of the experiment is to either improve ground station coordinates, the trajectory of the
satellite, or both. We will not only perform the measurements from just one location (read:
ground station or satellite) but instead involve more ground stations and satellites in a network.
Furthermore we will drop the assumption that the propagation speed equals to the speed of
light in vacuum, that perfectly synchronized clocks are used, or that the receiver or transmitter
or reflector stay at one place while we perform the measurement.

6.1.1 Satellite Laser Ranging

Since the invention of the laser in 1960 the possibility was created to send a coherent and high
intensive light pulse from a ground station into space. For the efficiency of hitting a satellite
with a light pulse the divergence of the laser should be as small as possible, in reality it depends
on the optical quality of the telescope. Also, the light pulse should be as short as possible, the
shorter the pulse the more unique it becomes, the better a returning pulse can be identified.
The detector in the receiving telescope is typically a photo-multiplicator tube or an avalanche
photo diode where a single photon can trigger the stop timer. The international organization
that collects SLR data for scientific purposes is called the ILRS, see also http://www.ilrs.org.
SLR tracking data may be collected up to the GPS orbit altitude by most tracking systems, an
attractive aspect of the technique is that we are only dealing with passive reflectors on satellites.
Reflectors for laser light can be build with high precision, the phase center of the reflector is
usually known within a few mm, and the divergence of the reflector can be make small. A
drawback is however that you need a manned tracking station and that you are depending on
weather conditions.

6.1.2 Doppler tracking

In figure 6.1 we assumed that there was a short pulse and that we measured the round trip
time of a pulse between transmitter and receiver. The end product was a range which finds its
use in navigation applications. There is one draw-back in the general concept, which is that a
short pulse is sent away, and that we are able to detect it. In the optical domain this is not a
problem, because the available bandwidth is very large, implementing a pulse is nothing more
than modulating a signal on top of a carrier. If the carrier has a frequency f and the modulated
signal a frequency g then sums and differences like f + g and f g will appear in the spectrum.
In the radio domain one faces the problem on bandwidth restrictions, hence pulse techniques do
not work in the radio spectrum.
For a radio tracking system one could even wonder whether it is necessary at all to modulate
a signal g with a substantial bandwidth on the carrier f . In Doppler tracking this is the case,

66
in principle nothing is modulated on the carrier f and just the Doppler effect of the carrier is
recorded at a receiver. To ensure some form of accuracy you need a high frequency stability of
both the oscillator at the transmitter and the receiver.
In order to track a satellite we can therefore make use of an easier principle to observe a
range: The satellite will have a velocity v and it travels along a direction n at a position indicated
by vector r relative to the receiver antenna (we call this the line of sight vector). The Doppler
effect seen by the receiver will be determined by the projection of v = n.v on the line of sight
r
vector r, hence vd = (n, ||r|| ).v. The Doppler effect seen by the receiver is therefore:
 vd  vd f
f + f = 1 + f f = f vd = .c (6.1)
c c f
where the inner product that leads to vd is computed such that vd is positive when the satellite
approaches the receiver. If we measure f then we observe directly vd , and this allows us the
construct a new type of measurement that results in the range change r(t0 , t1 ) between two
epochs t0 and t1 which are chosen somewhere in the acquisition interval.
Z t1 Z t1
f (t)
r(t0 , t1 ) = vd (t) dt = c dt (6.2)
t0 t0 f

Application of the Doppler range change effect is the basis of several tracking systems such as
DORIS but also GPS. In DORIS the Doppler effect is observed by a receiver in the satellite
while the beacons (the transmitters) are on the ground. With GPS the beacons are in the sky,
and the receivers are on the ground.

6.1.3 Global Positioning System

GPS is a technique that allows one to perform code and phase measurements to a space segment
that consists of approximately 30 satellites orbiting the Earth at around 20200 km. There are
several processing strategies, purpose is to improve the frequency and the phase definition of
a local oscillator which is used in the receiver clock. This first step is called the navigation
solution, it yields an approximate solution for the position and an approximate estimate for
the receiver clock bias. Next follow more advanced techniques to improve the quality of the
solution, and within the scope of the lectures we will discuss differencing techniques. In the
end we will mention two practical implementations which are realtime kinematic processing and
precise point positioning.

Receiver design
A GPS receiver is not like a normal radio receiver that can handle music or data, it works on
different frequencies in the L band. The original GPS design has the L1 frequency at 1575.42
MHz and the L2 frequency at 1227.60 MHz. The satellites are all on the same frequencies; spread
spectrum modulation is used to send out information to users on the ground. The modulation
scheme is rather advanced, but the essence is that the satellites modulate a so-called pseudo-
random noise (PRN) signal on the main frequency. For the PRN modulation you should know
that it consists of a clear access (C/A) code for the L1 frequency which is unique for each space
vehicle (S/V). There are also two variations of the PRN modulation, namely one that comes
down at a data-rate of 1 MHz and another at 10MHz. Figure 6.2 is meant to briefly summarize

67
Figure 6.2: Schematics of the GPS transmitter revealing the way the signal is modulated on
the L1 and the L2 carrier frequencies. Essential ingredients are the C/A code, the navigation
message and the P or the Y code, the difference between them is that the P code was publicly
known whereas the Y code is classified for civil applications. As of today all S/V emit the Y
code.

the main characteristics of GPS modulation. From this figure we can see that there is also
a navigation message data-stream that is superimposed on the PRN codes. These navigation
messages contain: precise timing information at the transmitting satellite, ephemeris of the
transmitting satellites, constellation almanac information and ionospheric model parameters.
The message is transmitted at a rate of 50 bits per second, it consists of 25 frames, each
frame has 5 subframes, that is made out of 10 words that have a length of 30 bits. Therefore
37500 bits need to go from the satellite to the user, at a rate of 50 bits per second this takes
12.5 minutes. After receipt the receiver knows where to find all satellites in the sky, and it will
collect the relevant frames of each S/V to retrieve the up to date parameters for the clocks and
the orbits. This is what happens at a cold start of the receiver, once completed the information
is stored in the receiver memory so that the stored information can used again when it has to
perform a warm start.
When you tune to GPS frequencies with an arbitrary radio then in first instance nothing
is heard. The signal is not picked up by a FM or AM demodulator circuit in your receiver,

68
instead, what you would need is a so-called BPSK (bi-phase phase shifted key) demodulator.
But even then there is an extra complication, namely that the signal is pushed to a 1 or 10
MHz bandwidth that most radio receivers can not handle, but also, due to the spreading of
transmitted energy the signal will appear for a user on the ground below the thermal noise level
of the receiver.
In order to tune to a specific satellite we need to demodulate the information. The only way
this can be done is to duplicate the C/A PRN code of a specified satellite and to cross correlate
this code with the incoming data. During the replication process we rely on the receiver clock
oscillator whose frequency is not precisely known, also, the phase offset of the receiver clock
is not exactly known, the situation is not really better than what your wrist watch performs.
In order to navigate with any accuracy (say 3 meter) the receiver clock has to be modelled to
better than 10 nanoseconds.
There are two main effects here that the receiver needs to handle. First is that we dont
know at a hard start where the satellite is, so the Doppler shift of the transmitted information
has to be guessed. Second, we do not really know what the phase offset of the transmitted C/A
code is. What we do know however is that the C/A code repeats itself in 1 milli second, for the
P/Y codes this is another story, in that case the PRN repeat time is of the order of 20 weeks.
For this reason the GPS receiver initially tries different C/A code phase offsets and guessed
frequency variations for each GPS S/V until it finds a match in the code and frequency domain.
Once it has a lock on one satellite it will download the full almanac and it will try at the same
time to demodulate information from other GPS S/Vs. When this process is completed we can
go to the next step, which is to use to C/A code phase information and the GPS S/V orbit and
clock information to carry out a navigation solution.

The navigation solution

As was explained before, PRN code information in the GPS receiver is nothing more than a
phase difference between what the satellite transmitted and what the receiver clock assumed.
There is also an ambiguity in the process because the C/A PRN code repeats itself every 1
millisecond which is 300km in length. But, there is C/A PRN information from more than
one satellite so that the combination of several satellite codes results in a situation where we
approximately know where the receiver is on the Earths surface.
The user can help a GPS receiver here, help in the sense that it speeds up the initial C/A
code untangling process where we dont exactly know how many code multiples of 1 msec exist
between the receiver and the satellite antenna. But again, you can also leave it up to most
receivers nowadays, they will be able to find the most likely spot on Earth where the 1 msec
multiples to the space segment match. After this process has completed we fixed the receiver
clock offset to within 1 msec. Is that good enough to navigate around, no, it is not, so for this
you need to perform some mathematical modeling where you combine information from at least
four different GPS S/Vs above your local horizon.
Any C/A code observation for which the receiver found the 1 msec code multiples relative to
the space segment can be seen as a so-called pseudo-range measurement. This measurement
is just like any range measurement, except that there is a bias that is mostly caused by the
uncertainty of the receiver clock which was already modelled to within 1 msec in the C/A code
untangling procedure that we earlier described. Therefore:
(tr , ts ) = c(tr ts ) + c.(tr ts ) +  (6.3)

69
where tr is the receiver time, ts is the satellite time, c(tr ts ) is the geometric distance between
satellite and receiver where c is the speed of light. The term c.(tr ts ) specifies a range
effect due to the receiver and satellite clock error. The latter are small because GPS S/Vs
are equipped with Rubidium and Cesium frequency standards that are continuously monitored
from the ground. Finally  is everything we conveniently ignore at this point because our first
interest is to improve the accuracy of the receiver clock.
If we combine four pseudo ranges, and if we assume that the receiver location is roughly
known (we already know it to within 300 km because of the C/A code untangling procedure)
then it is also possible to linearize equation (6.3):

(tA par /Xr pA A

r) r /Yr pr /Zr 1 Xr
(tB pB
r )
B B
r /Xr pr /Yr pr /Zr 1 Yr


C =
C C C
(tr ) pr /Xr pr /Yr pr /Zr 1 Zr (6.4)
(tD r ) pD D D
r /Xr pr /Yr pr /Zr 1 ctr

where the partial derivatives in the matrix are computed at the linearization point, which is
our guessed initial position of the GPS receiver. After solving this system of equations we
obtain improvements of the receiver coordinates, and an estimation of the setting to apply to
the receiver clock. New partial derivatives may be computed and the procedure can be repeated
with updated coordinates and clock offsets for the next batch of C/A codes produced by the
demodulation circuit. Put this information in a Kalman filter (will be discussed later in these
notes) and you are able to produce velocity and course heading of a moving vehicle. For 99% of
the applications of GPS this is sufficient, there is a small number of users who want to obtain
better navigation and clock information.

Advanced GPS measurements

Scientific GPS receivers are not only able to receive the C/A codes and perform a navigation
solution, but they also provide the user with carrier phase information and ionospheric delays
between the Y-codes at L1 and L2 which can be cross correlated during the demodulation
process. In the end the receiver is able to integrate the Doppler curves, cf. the Doppler tracking
section 6.1.2, because the instantaneous Doppler frequency of a GPS S/V is produced by the
frequency tracking loop in the demodulation circuit. To summarize the discussion: a scientific
GPS receiver is able to provide the C/A code on L1 , a carrier phase on L1 and a differential
code measurement between L1 and L2 . By code squaring most receivers can strip the Y codes
from the signal, in the end this results in the carrier phase measurements on both L1 and L2 .

Single, double and triple differencing

One of the main difficulties with the navigation solution is that the receiver clock can not
be modelled better than the behavior of the satellite clocks. The quality of these clocks is
impressive, but it does not impress a scientist who insists on sub centimeter knowledge of GPS
antenna phase centers. A some point in time in the past GPS was deliberately distorted by DOD
to prevent that users could compute their positions and clocks in real time with high accuracy.
This is known as selective availability which was turned off after a presidential order of Bill
Clinton in May 2000. Most scientific users are not interested in real time GPS navigation, and
are perfectly happen when receivers are connected in a network so that data can be exchanged

70
(by independent radio communication techniques either on ground, or via a separate satellite
link) and processed later or in real time. It should be obvious that we are able to further improve
the quality of the GPS solutions, the easiest example is to implement a differencing technique.
If there is a remnant receiver clock error tP then all code and carrier phase data from
that receiver will be biased by a range effect as large as c.tP . Suppose now that there are
two independent satellites A and B, in this case the difference (tP , tA ) (tP , tB ) will be
not be affected any longer by the receiver clock bias tP . The quantity SD(tP , tA , tB ) =
(tP , tA )(tP , tB ) then said to be a single-difference formed by receiver P relative to satellites
A and B. Single differences can also be formed by a satellite to two ground receivers, in this case
we get the observation SD(tP , tQ , tA ) = (tP , tA ) (tQ , tA ). Single differences by themselves
are still affected by either two satellite clock errors, or two receiver clock errors and this can be
improved by processing the single differences one step further to end up with double differences.
The quantity DD(tP , tQ , tA , tB ) = SD(tP , tA , tB )SD(tQ , tA , tB ) will be free of any clock error,
but it requires one to establish an independent communication path between receiver A and B
on the ground. To summarize, double difference are free of all possible receiver and satellite
clock errors that still remain after the navigation solution.
However, we did not yet explain you how (tr , ts ) was provided by the GPS receiver. The
easiest way would be to use the C/A pseudo range data, but this is not very accurate because
of the way C/A code data is demodulated by the receiver. C/A codes repeat themselves every 1
millisecond, they consist of 1024 code chips that at roughly one microsecond long in time, and
the digital code correlator in your GPS receiver will never be able to perform the code phase
measurement better than typically 1/100 to 1/1000 of the length of the code chip. Somewhere
between 1 and 10 nanoseconds is a very realistic estimate for the C/A code phase measurement.
A method to improve this is to use the phase measurements relative to the carrier, which we
also get from the GPS receiver, but at a price which is that the integrated Doppler tracking of
the carrier frequency resulted in a range change. This a range that is relative to an assumed
reference epoch t0 where we started the L1 and L2 carrier phase tracking loop integrators in the
GPS receiver. The key point is here that t0 is arbitrary, we can only start to integrate when the
satellite appears above the horizon.
If double difference information can be formed from the C/A code measurements then also
the same quantity can be formed from carrier phase tracking data. The change in time of a double
difference quantity can now be replicated when all carrier phase data is used, and this results in a
so-called triple difference T D(t1 , t0 , A, B, P, Q) = DD(t1 , A, B, P, Q) DD(t0 , A, B, P, Q) where
we use a short hand notation for the involved double differences. The benefit of triple differences
is that they are independent of the carrier phase start values, and that they are more accurate
because they can be formed from the carrier phase information which has a range error of the
order of a mm while pseudo ranges are no better than 30 to 300 cm one sigma depending on the
technology of the receiver. After single, double and triple differencing techniques are applied we
can construct a network solution where the relative coordinates and clocks between the receivers
is reconstructed to high accuracy, typically better than 10 millimeter or 33 picoseconds. With
such information we can do geophysical research such as to measure velocities between tectonic
plates.

71
RTK: real time kinematic processing
If differencing techniques are applied between a stationary GPS base station A and a roving
vehicle B and if all demodulated information is sent in real time from station A to B then the
user at B has the ability to navigate with a relative accuracy that is significantly greater than
that he/she would have performed a stand-alone navigation solution. Typically RTK is used
for civil engineering construction work, or it may be used for landing an aircraft without the
availability of DME or ILS at an airport, or it may be used to guide a military jet at night
towards the landingstip on an aircraft carrier.

PPP: precise point positioning

This technique is similar to RTK except that the internet is used to distribute for instance
satellite clock corrections and other information such as carrier phase offsets so that the user can
independently benefit from the network processing. The Fugro company sells PPP information
to maritime users, for this you have to buy a separate receiver to demodulate a BPSK modulated
signal from an Inmarsat satellite.

6.2 Corrections
6.2.1 Light time effect
Range and velocity measurements need a correction for the light time effect when the receiver
or the transmitter or the reflector move relative to one another. Let us consider the simplest
example in SLR where a laser is shooting at a satellite at an altitude of 500 km passing overhead
at a speed of 8 km/s. The round trip time for laser light (disregarding any other measurement
errors) is 3.33 msec and during this short time the satellite (and the ground station) will move
because of the 8 km/s and a maximum of 464 m/s due to Earth rotation. At 8 km/s with 3.33
msec light time this means that the satellite has moved over a distance of 26.6 meter which is
significant for the problem, in particular when the motion projects on the line of sight along
which the measurement is performed. The remedy is to use a-priori models that predict the
configuration, and to apply the light time effect in these models to find out where the reflection
occurred, and next to correct the SLR measurements to account for the light time effect. The
first-order correction is usually sufficient, there is not really a need for iteration with this problem
unless the observed ranges become large, for instance in interplanetary applications. During the
class we will ask you to compute the light time correction in matlab for a given configuration.

6.2.2 Refraction
Refraction is caused by the fact that an electromagnetic wave (either light or radar) has to
travel through a medium where the speed of propagation v is less than the speed of light c. The
refraction index n is nothing more than c/v and has in most cases a value greater than 1. If the
refractive index of the medium is known, and if we travel from a transmitter A to a receiver B
via a refractive medium we will find a correction for the range s which is caused of the physical
properties of that medium. The procedure is as follows:

s = c.t = v.(t + t)

72
where t is the theoretical travel time t in vacuum, and where t+t is measured by an instrument.
Due to refraction we get: c
t + t = t = n.t
v
By rearranging terms in this equation we get:

(n 1)t = t

so that:
s = ct
The last step is to integrate over all s along the path from transmitter to receiver, this results
in the Fermat1 integral: Z B
s = (n(s) 1) ds (6.5)
A
which gives us the range correction s. This integral should be evaluated along the curved path
between A and B where curvature of the light ray may be described by Snells law. This law
states that, if a wave hits an interface between medium 1 with refractive index n1 with an angle
1 with respect to the normal vector on that interface, then the wave will continue with normal
angle 2 in the second medium where the refractive index is n2 . In that case Snells law states
that n1 sin 1 = n2 sin 2 see also figure 6.3
The consequence of refraction will be discussed for three well known refractive effects in the
atmosphere where we distinguish between the dry tropospheric effect as a result of gas molecules
like O2 N2 and other gasses, the wet tropospheric effect due to H2 O in gaseous form (thus not
in condensed form like rain etc) and the ionospheric effect which is a result of free electrons in
the thermosphere caused by the Sun (and the interaction between charged particles originating
from the Sun including the interaction with the Earths magnetosphere).
In reality it is important to distinguish between the group and the phase speed of a wave
of which we said that the it traveled at speed v. The phase speed of a wave relates to the carrier,
if nothing is modulated on the wave then this is the only propagation speed relevant for the
calculation. However, if information is superimposed on the carrier by means of a modulation
technique, then we also have a group speed which counts for the information that is modulated
onto the carrier. Also, a medium can be dispersive, or it can not, meaning that the propagation
speed (either group or phase) depends on the frequency of the wave. In case of a dispersive
medium we will find different values for the group and the phase speed of the wave. The
consequence of dispersive and non-dispersive refraction will be discussed in the following where
we discuss two cases, namely a laser and a radar altimeter system that measure their altitude
above the Earths surface.

Dry tropospheric effect

If a microwave radar signal from an altimeter in space travels to the sea surface beneath the
satellite then the total range correction caused by the ionosphere and troposphere becomes
relevant. In both media the refractive index differs from the value of 1. The dry tropospheric
effect is not dispersive for microwave frequencies, and this property extends to the optical domain
where dispersion caused by dry gasses is small so that it is usually ignored. As stated before, the
1
Pierre de Fermat, lawyer, statesman and mathematician 1601-1665

73
Figure 6.3: Snellius law, also known as Snels law.

74
dry tropospheric correction describes the delay of a radar signal caused by the presence of gasses.
The total effect counted from the mean sea surface vertically to space amounts to approximately
2.3 meter and mainly depends on air pressure at MSL. If a range measurement is performed
at a certain height above MSL, then station height should be taken into account. Meteorologic
models that provide sea level pressure (or geopotential height) are normally employed. The
accuracy of the correction depends on the quality of the supplied air pressure data. As a rule
of thumb meteorological models provide air pressure to within 3 mbar standard deviation on
a global basis which should be compared to the nominal value of 1013.6 mbar of the standard
atmosphere. This means that the relative accuracy of the dry tropospheric correction is no more
than 0.3% which translates to 6.9 mm. Slant ranges through the troposphere should take into
account the geometric effect of a longer path through the atmosphere, it is up to the reader
to verify that the lengthening is proportional to 1/cos Z where Z is the zenith angle. Better
functions such as the Vienna mapping function have been developed over the years to account
for slant ranges, the goal is then to estimate the vertical delay as one parameter and to rely on
the mapping function for other values of Z, for details see [63].

Wet tropospheric effect

Laser light is not affected by the wet troposphere, but all adio frequencies are affected by the
refraction of water gas. The wet tropospheric effect is related by the presence of humidity and
this should not be confused with condensed water in the atmosphere which scatters and therefore
attenuates radio signals. For a number of reasons it is a significant effort to accurately compute
the wet tropospheric correction of radio signals in the Earths atmosphere. The nominal delay
for the wet troposphere in the tropics can be as large as 40 cm and at higher latitudes it gets
down to approximately 10 cm, for arctic regions it is negligible. Oftentimes meteorologic models
are not accurate enough, if they were used to compute a correction then 60% is a very optimistic
estimate for relative accuracy. In practice this means that more than 5 cm noise easily remains
thereby introducing one of the biggest difficulties for instance in designing an accurate altimeter
system whereby the end user (oceanographers) asks for less than 5 cm noise in the sea level
reconstruction. The remedy is to install a water vapor radiometer (WVR) on the spacecraft to
measure the brightness temperatures of the Earth at two or three lines in the electromagnetic
spectrum near the water vapor absorption line at 22 GHz. Some altimeter systems, such as
GEOSAT (19851990), did not carry a WVR and external oblique space-borne radiometric data
had to be used to provide a correction. With the aid of WVR data on several altimeter systems
since the 90s the wet tropospheric correction can usually be modeled to within 2 cm or so. For
GPS other techniques should be used to correct the wet tropospheric effect.

Ionosphere
Ionospheric refraction is dispersive (frequency dependent) and the effect can be estimated by
measuring ranges at different frequencies. For lower frequencies the ionospheric effect becomes
more pronounced. Below 30 MHz the ionosphere is a reflective layer, and this enables HF radio
communication behind the local horizon. Phenomena such as fading on the MW band during
the night on your radio (526.5 to 1606.5 kHz in Europe) are caused by the Earths ionosphere.
During the day ionization occurs because of sunlight and the consequence is that MW signals
will not propagate too far, but during the night you may be able to listen to stations which are

75
several thousand kilometer away from you. The MW signals bounce from the transmitter to
the ionosphere, and bounce again on the surface, numerous hops allow you to bridge incredible
distances, but, the effectiveness of the propagation depends on the state of the ionosphere for
MW signals. Between MW frequencies and 30 MHz, or below MW frequencies other phenomena
play a role but similar propagation properties are found. Above 30 MHz radio waves usually
pass through the ionosphere where refraction becomes dispersive.
For frequencies in the electromagnetic spectrum near 13 GHz the ionospheric correction is
usually less than 30 cm so that most radar altimeter satellites had to rely on the presence of
global ionospheric models to compute the correction. The accuracy of the global ionospheric
models is probably of the order of 90% meaning that 3 cm noise is introduced by this way of cor-
recting the data. With the advent of dual-frequency altimeter systems such as TOPEX/Poseidon
information became available allowing to remove the first-order ionospheric delay from the ob-
served ranges because of the dispersive nature of the ionosphere. Essentially we can obtain the
correction from:

r1 = r0 + 2 (6.6)
f1

r2 = r0 + 2 (6.7)
f2
where r1 and r2 are measured at frequencies f1 and f2 and where we intend to obtain the
true range r0 and the ionospheric parameter . Usually there is more short term noise in r0
compared to the short term noise in r1 and r2 . For this reason moving average filters can assist to
reduce short periodic noise in the dual-frequency measurements of the ionospheric delay. After
smoothing the correction the ionospheric error is typically less than 1 cm for dual frequency
altimeter systems.

6.2.3 Multipath
Under ideal conditions we want to observe ranges along the direct line of sight between trans-
mitter and receiver, we dont want that additional reflections occur, and that those reflections
affect the distance measurement. In optics you can not make the mistake because the mea-
surement does occur at very high frequencies where the divergence angle is only controlled by
the aperture of a sending and a receiving telescope. But in the radio domain frequencies are
lower, and divergence is usually larger. Because of the latter is not unimaginable that a radio
signal first bounces on a reflector, and that it is next picked up by the receiver, or that both
phenomena take place at the same time so that the receiver sees a mixture of the direct signal
and the reflection.
Before the era of cable TV ghost images frequently appeared on television sets when signals
were weak. When an aircraft crossed the line of sight between the receiving antenna and the
transmitter a wobbly mirror of the direct signal appeared caused by the signal reflecting on the
aircraft. Multipath is only avoidable by removing all reflectors between the transmitter and the
receiver, we see it with GPS to some extend despite the fact that countermeasures are taken
to suppress multipath. In GPS circular polarization is used, the transmitter sends a righthand
polarized signal and the receiver antenna is mostly sensitive to the right hand polarization. A
reflection results in a weaker signal, in addition the sense of polarization changes from right-
hand to left-hand, so that the GPS receiver antenna wouldnt pick it up. But the reality is that
antennas are not perfect and that some of the reflected signal enters the GPS receiver.

76
 Multipath is more of an issue with GPS code measurements than that it is with the phase
measurements that have a short wavelength (19 and 24 cm respectively with GPS), compared
to the code measurements that come with a chip length of roughly 10 microseconds which
is 300 meter. Multipath depends on the wavelength, and long waves are more significantly
affected by reflections than short waves. Phase measurements can be done with an accuracy
of approximately 0.1% of the wavelength, the code phase discriminator noise in a typical GPS
receiver is approximately 30 cm, for carrier phase the phase discriminator noise is typically 1
mm.
A well chosen observation sight in therefore above the local ground (which does leave reflec-
tions), and sufficiently far away from anything that may cause a reflection. Still, multipath can
occur within the antenna, and some designs are better than others in suppressing the effect.
Well know is the Dorne-Margolin design that has concentric rings around a phase center.
During a data analysis multipath becomes visible when the residuals of either code or carrier
phase measurements are plotted in an antenna phase diagram. (Simply: determine the azimuth
and zenith of your raw measurements and start averaging all residuals in the data analysis in
this frame). Any appearance of a concentric ring shaped pattern is indicative for the presence
of multipath, the cause of the problem should be investigated. Alternatively, antenna phase
correction maps may be constructed to suppress multipath effects. Some manufacturers provide
antenna correction maps that were constructed in an anechoic chamber.

6.3 Applications
6.3.1 Satellite altimetry
Altimetric measurements
The measurement principle is, conceptually seen, an observation of the shortest distance from
the radar antenna phase center to the sea surface beneath the spacecraft (S/C). Contrary to
popular belief there is no such thing as the perfect pulse-radar, instead modern altimeters are
based upon a frequency modulation (FM) technique where a linear chirp signal with a dynamic
range of 300 MHz is modulated on a 13 GHz carrier, see also figure 6.4. The carrier and
modulation frequencies are just mentioned as examples and differ somewhat from the actual
frequencies used for the TOPEX/Poseidon Ku-band altimeter2 . After receiving the chirp signal
it is multiplied by the transmitted signal which allows to derive the frequency difference being a
measure of distance. Certain ambiguities may occur which are in general avoided by choosing a
proper modulation scheme and minimizing the altitude variations with respect to the sea surface.
The difference signal labeled T-R in figure 6.4 is then provided to a Fast Fourier Transform
processor returning the raw waveform samples. From this figure it is obvious that the inverse
Fourier transform of the T-R signal is equivalent to a phase (or distance) measurement of
two saw-tooth signals and that the FFT processor will simply return a histogram of observed
ranges. These radar waveform samples should be seen as the response of a sea surface equipped
with wind waves to a short (but not infinitely short) radar pulse.
Normally far too many raw waveform samples are generated so that statistical filtering
2
The T/P altimeter system operated between 1992 and 2006, it was launched by NASA and it was at its time
a revolutionary instrument that provided unique insights in ocean tides, mesoscale variability, sea level rise, and
Kelvin and Rossby wave propagation in the oceans

77
 Figure 6.4: Altimeter schematics based upon linear frequency modulation.

methods based upon alpha-beta trackers or Kalman filters are applied to smooth and compress
the data stream. This procedure is also carried out on board the S/C and is optionally controlled
from ground for certain altimeter systems. For the T/P S/C (see the JGR oceans special issue
of December 1994) one obtains 10 Hz ocean-mode waveform data which include range estimates.
For ERS-1 and ERS-2, two altimeter launched by the European space agency in the 90s, there
are two programmable modes, one for flat ocean surfaces and another for rugged ice surfaces.
The altimeter ice-mode is designed around the philosophy of measuring a wider domain of
distances with decreased accuracies. This is accomplished by reducing the saw-tooth frequency
range and relaxation of leading edge alignment criteria during the statistical processing of the
raw waveform samples.
The linear FM radar technique described above has the main advantage that power-hungry
pulse radar methods can be avoided and that low-power solid-state electronics may be applied
reducing the cost of implementing the radar altimeter. Clearly the radar waveform data are
nothing more than a distribution of the reflected ranges in defined time slots. The typical shape
of the radar waveforms is of course determined by the dielectric properties of the illuminated
surface, the antenna divergence and the off-nadir pointing angle of the altimeter. This illumi-
nated sea surface, or radar footprint, is roughly 10 km in diameter depending on the state of
the sea surface.

Radar correction algorithms

There are several radar corrections which are specific to the altimeter, they are however essential
for obtaining high precision altimeter measurements and determine the system accuracy that

78
can be achieved. The most frequently applied corrections concern the altimeter EM-bias, iono-
spheric and tropospheric delays, and the inverse barometer effect. All corrections that deal with
refraction was already discussed, what remains for satellite altimetry are the electromagnetic
bias and the inverse barometer effect.

EM-Bias
The sea surface is a well known reflector causing radar waveform samples to be determined by
wind waves in the radar footprint. A measure for the roughness of the surface is the so-called
significant wave height (SWH) parameter which can be derived by fitting a well chosen function
through the waveform samples, cf. [25]. SWH varies between approximately 0.5 meters up to
20 meters with a global average of about 2 to 3 meters. Low SWH values usually indicate
reflections of flat surfaces such as ice. Editing the altimeter data is crucial before you apply it in
any further analysis. Extreme SWH values usually indicate storm conditions so that the validity
of this correction will become a problem. The EM-bias correction is the result of the asymmetric
shape of the sea surface since more radar signal is reflected from wind wave troughs than crests.
It inherently leads to an electromagnetic bias or EM-bias since the measured surface will appear
somewhat different than the real surface. Some remarks:

The EM-bias correction should estimated from the SWH parameter prior to using the
altimeter data. More sophisticated algorithms for the EM-bias correction incorporate
knowledge about the wind speed (U ) at the sea surface. The scalar wind speed U is
estimated from other characteristics of the radar waveform samples, see also the discussion
on the Brownian reflection model cf. [3].

Typically the EM-bias correction is of the order of 3% with respect to the SWH with an
uncertainty of the order of 1%. In early altimeter research the EM-bias effect was only
estimated as a scale factor relative to the measured SWH, a higher regression constant
(7%) was found for the SEASAT altimeter. In modern research the effect also involves a
wind speed regression constant.

The electromagnetic-bias is caused by an asymmetric reflection at the ocean surface, the

effect should not be confused with the sea-state bias correction applied to altimeters since
the latter also includes the waveform tracker biases. Thus: SSB = EM-bias + tracker bias.
The SSB effect is usually obtained from a repeat pass analysis of the altimeter3 .

The sea state bias of the altimeter is a fundamental limitation of satellite altimetry, the
consequence is that in-situ altimeter height readings and not better than approximately 1
cm over the oceans. We can only improve the analysis by involving more altimeter data
in a spatial regression analysis.

Inverse Barometer Correction

Apart from its role in computing the dry tropospheric range correction, air pressure will affect
the sea level which responds as if it were an inverse barometer. In this case we will see that
3
The tracker is the algorithm in the satellite that detects the leading flank of the returning waveform which
is used for the radar range estimation, the re-tracker is the algorithm that usually runs on the ground, input are
the waveforms from the altimeter

79
there is a linear relation of almost -1 cm per mbar; the minus sign tells that the sea level
is depressed by 1 cm when air pressure is increased by 1 mbar, hence its name: the inverse
barometer mechanism. The practical way of dealing with the problem is to use a meteorologic
sea level pressure dataset so that the inverse barometer correction itself may be computed to
better than 3 cm. Nevertheless some remarks should be made:

In reality the inverse barometer correction is more complicated than a simple linear re-
sponse; reason is that the barometric effect is forcing the ocean via a pressure gradient
and a wind stress curl. Accelerations on water parcels due to the Coriolis effect, friction
etc should be taken into account when you solve the Navier Stokes equations,

It turns out that the inverse barometric response is not very effective on time scales under a
day or so, to model this there are better approaches that take into account high frequency
aliasing effects,

On the 1 mbar level tidal signals exist in the atmosphere and one should find out whether
the ocean tide model is in agreement with the pressure models being applied,

In the tropics the natural variation in air pressure is small compared to other regions on
Earth and statistical analysis of altimeter data, ie. comparison of air pressure variations
against height variations of the sea surface in the tropics and also in western boundary
regions, has shown that the -1 cm per mbar response is not per se valid, cf. [33].

Altimeter timing bias

The timing bias correction originates from the era of Skylab (1974) where the altimeter was
activated by manual control by an astronaut within a space station. Since the human delay
in turning on an instrument is probably of the order of 0.5 second it means that altimeter
measurements were recorded with the wrong time tag. So the measurement itself was correct,
but the time stamp in the dataset with altimeter records was shifted by 500 msec. Initially
the time tag problem was ignored and interpretation of the Skylab data soon suggested that
existing maps of the Earths geoid had substantial errors in excess of 10 to 20 meter. The
latter was certainly not the case, instead, the error was man-made. Later it was realized that
the altimeter range error was correlated with the vertical velocity of the Skylab space station
which varies within -25 to +25 meters per second. The effect of the altimeter timing bias is
as straightforward as multiplying the vertical speed of the s/c above the sea surface with the
timing bias. Consequently in order to obtain acceptable values of less than 1 cm it is required
to get the time tags to within 500 sec.
It turns out that there are no fail-safe engineering solutions to circumvent altimeter timing
error problems other than to calibrate the S/C clock before launch and to continuously monitor
the S/C clock during the flight via a communications channel with its own internal delays. A
practical way to solve the altimeter timing bias problem was suggested by [36], the method is
based upon estimating a characteristic lemniscate function that will show up in the altimeter
profile. Nevertheless ERS-1 still exhibits time tag variations at the 1 msec level according
to (Scharroo, private communications) which corresponds to 2.5 cm mostly at two cycles per
revolution. A better approach would be to rely on GPS timing control in the S/C which is
normally better than a few nanosec. Remnant time tag effects could still remain, oftentimes

80
 Figure 6.5: Skylab space station, image source: NASA

they are caused by correctable software/hardware anomalies in the processing of the telemetry
data on the ground.

Surface slope corrections

Radar altimeter distances perpendicular to an inclined terrain introduce a range effect of the
order of 2 where is the surface gradient and the distance. Since 106 meter should
be smaller than 104 radians which corresponds to a terrain gradient of 10 meter over 100 km.
Normally such gradients are not encountered over the ocean unless the altimeter is used to map
the geoid over subduction zones. The terrain gradient effect is more important for altimetry
over land ice. The effect is very significant over the margins of the ice caps on Greenland where
the terrain slope correction is clearly a problem for a meaningful interpretation of the data.
In this case crossover processing or a digital elevation model help to suppress the terrain
gradient effect. But the other design aspect is to build an altimeter that can operate such that
slant ranging is enabled over the ice sheet margin. This is implemented in the SIRAL instrument
on the CryoSat-2 altimeter which comes with two interferometric modes, one SAR mode is for
high resolution mapping over sea ice, the other SARin mode is for mapping over sloped terrain
whereby two altimeter antennas are used.

Nominal orbits
Two important aspects playing a role in the choice of a nominal orbit are the minimization of
the altitude variations and a global coverage of the ground track pattern. The first requirement

81
is more or less a consequence of the altimeter design itself, the second is more a practical desire
for an Earth observing mission. Additional requirements can be for instance a full sun orbit
(as realized for the SEASAT mission) or a Sun-synchronous orbit (as with ERS-1) whereby the
local solar time at the sub-satellite point is always at the same value facilitating the operation of
Earth sensing instruments. For certain oceanographic applications a trajectory may be chosen
such that the ground track at the Earths surface repeats itself after a certain period which is
also known as the repeat cycle.
The need for frozen orbits arises from the minimization of altitude variations. The relation
between eccentricity and argument of perigee in low eccentricity orbits is given by [10] who found
that precessing, liberating and frozen conditions may occur depending on the choice of the initial
orbital elements and moreover the ratio between the odd zonal harmonics and J2 of the Earths
gravitational field. The theory was implemented in the design of the TOPEX/Poseidon orbit
where the mean value of the argument of perigee is fixed at = 270 by adopting a specified
mean eccentricity at the (0, C/k) point in Cooks (u, v) diagram of non-singular elements. Un-
fortunately the frozen orbit is unstable and requires periodic maintenance of about once per
month.
In these lecture notes we will frequently refer to the results obtained by the TOPEX/Poseidon
altimeter satellite which was placed in a frozen repeat orbit at 1330 km altitude. The orbit itself
is not synchronized relative to lunar or solar motions which is extremely helpful for analyzing the
ocean tide signal, see also section 16. The repeat conditions of the T/P orbit are such that the
ground track overlaps in 127 revolutions or 10 nodal days which corresponds to 9.9156 wall-clock
days. The ground track spacing between parallel altimeter tracks is therefore 360/127 = 2.8
where the ground track will reach the extreme latitudes of 66 at the Earths surface.

Precise orbit determination

Any altimeter satellite places extreme requirements on the quality of the orbital ephemeris.
The goal is to compute the position of the center of mass of the spacecraft to better than 3
centimeters. This task turns out to be a very difficult orbit determination problem that could
not result in the desired radial orbit accuracy for many years. Initially the radial position error
of the SEASAT and GEOS-3 altimeter satellites was typically 1.5 meter manifesting itself as
tracks in altimetric surfaces which are clearly identified as 1 cycle per revolution orbit errors.
It turned out that these radial orbit excursions were mainly caused by the limited accuracy of
then existing gravity models.
However altimetry as a technique is not useless because of poor orbits and considerable
effort went into the design of efficient processing techniques to eliminate the radial orbit effect
from the data. Collinear track differences are mostly insensitive to gravity modeling errors
and rather efficient adjustment techniques enable to remove radial trend functions between
overlapping altimeter tracks. Many papers have shown that such processing schemes result in
realistic estimates of the oceanic mesoscale variability which normally doesnt exceed the level
of approximately 20 to 30 centimeters.
Other processing schemes are based on a minimization of cross-over differences which are
obtained as the sea surface height difference measured by the altimeter at intersecting ground
tracks. In an attempt to reduce the orbit error, linear trend functions are estimated from
relatively short and intersecting orbits. Another possibility is to represent the radial orbit error
as a Fourier series for a continuous altimeter arc spanning a full repeat cycle. A summary on

82
the efficiency of such minimization procedures is discussed in more detail in [55] where the null-
space problem in relation to cross-over adjustment is discussed. In these lecture notes you can
switch to section 8.8 where compatibility conditions of rank deficient linear sets of equations are
discussed.
In fact the problem is rather similar to determining absolute heights from leveling networks
where the measurements are always provided as height differences between two stations. The
mathematical solution is to apply at least one constraint, known as a datum, which determines
the height offset of the network. However, in =case of cross-over minimizations the datum
problem is ill-posed and fully depends on the assumption of the orbit error trend functions. Due
to rank deficiency only a partial recovery of the orbit error function is feasible which will obscure
the long-wavelength behavior of the sea surface modeled from altimeter data.
There are several reasons for not applying collinear or cross-over adjustment techniques
in contemporary precision orbit determination schemes. A first reason is that we have seen
significant advances in modeling the Earths gravitational field. The older gravity models, such
as GEM10b, cf. [34], were simply not adequate in describing the rich spectrum of perturbations
of an orbiter at 800 km height such as ERS-1, GEOSAT and SEASAT. The Joint Gravity
Model 2, (JGM-2, also named after the late James G. Marsh who was one of the early pioneers
in precision orbit determination and satellite altimetry), is now complete till degree and order
70, cf. [51]. A second reason is in the design of the TOPEX/Poseidon itself whose altitude is
1330 km which inherently dampens out gravity modeling errors. Improved tracking systems such
as DORIS (a French Doppler tracking system on TOPEX/Poseidon and several other altimeter
satellites) and a space-borne GPS receiver, cf. [19] have completely changed the situation in
the beginning of the 90s. The result is that the orbit of TOPEX/Poseidon can be modeled
to less than 2.5 cm rms which has completely revolutionized the processing strategy, and more
importantly, our understanding and interpretation of altimeter data. More recent re-processing
of all altimeter data since the early 90s has shown that the orbit error can be reduced further
to approximately 15mm in error.
Nevertheless there are still a number of open problems that could stimulate future research.
First of all we think that there is still valuable information in the existing backlog of altimeter
data from GEOS-3 and onwards where orbits and altimeter corrections may require reprocessing
which ultimately may help to better understand the behavior of the Earths oceans and in par-
ticular the modeling of inter-decadal fluctuations of the sea surface. In our opinion a worthwhile
experiment would be to recompute all the existing altimeter orbits preferably in a reference
frame consistent with TOPEX/Poseidon and moreover to re-evaluate all models required for
reprocessing the altimeter data records.

6.3.2 Very long baseline interferometry

VLBI is a technique to map the intensity of astronomic radio sources with the help of dedicated
antennas and receivers. The radio sources are natural and consist for instance of quasars which
are thought to be massive black holes pulling in material that is accelerated to very high speeds.
Radiation is then emitted in the form of X-rays, ultraviolet and visible light, but also in the
form of synchrotron radiation which can be observed by a radio telescope. Other radio sources
are for instance neutron stars which as thought to be remnants of a massive star that collapsed,
the neutron star rotates and it can be observed with radio telescopes.
On Earth we are able to detect radio emissions from natural sources with one or more tele-

83
 Figure 6.6: Left: Westerbork synthetic radio telescope, Right: a receiver

scopes. Figure 6.6 shows an array of antennas at the Westerbork observatory in the Netherlands.
The combination of antennas allows one to simulate a very large dish which helps to increase
the angular resolution of the instrument. Angular resolution is a property of any device that
receives electromagnetic radiation, regardless whether it is a radio telescope, your eye, a UHF
antenna or a satellite TV dish. The general rule is that the shorter the wavelength the better
the resolution, but also, the larger the aperture the better the resolution will be and visa versa.
The best instrument has a large aperture and detects small wavelengths. Why is this the case?
Figure 6.7 shows the relation between the wavelength of electromagnetic radiation in
relation to an antenna whose aperture is greater than . Imagine that an isotropic radiator is
placed on both sides of the aperture with dimension D. Constructive interference will now occur,
that is, there will be a maximum in a plane perpendicular to the pointing axis of the antenna
where both isotropic radiators are in phase. If we place two coherent radiators at both sides of
the aperture then extinction will occur for an imaginary wavefront that is /2 out of phase at
the top of the aperture and in-phase at the bottom because of interference. This happens also
for an imaginary wavefront that is in-phase at the top and /2 out of phase at the bottom, see
also figure 6.7. The angle indicated in figure 6.7 follows from the ratio of over the aperture
D of the antenna, a good approximation is that = /D, the better approximation is that

= arctan( D ) This approximation is only valid when the antenna aperture is substantially
larger than the wavelength .
For larger wavelengths, and hence lower frequencies, antennas are usually constructed out
of dipoles so that the the /D approximation is no longer valid. The dipole itself is a /2 bar
of metal that comes with an excitation point the center where the dipole is cut in two parts,
it is a isotropic radiator with a reduced sensitivity along the dipole axis, the radiation pattern
for a vertical positioned dipole is shown in figure 6.8 on the left. The right part of figure 6.8
is an example of an array of dipoles placed on a support structure. This method significantly
increases the angular resolution of the UHF antenna, the pointing ability becomes a couple of
degrees depending on the number of dipoles in the Yagi array and the wavelength that is used.
For the angular resolution of your eye the wavelength of light is divided by the aperture
which is the diameter of your iris. The angular resolution of an unaided eye is therefore 500nm
divided by 8mm (healthy young eyes can do this), which is 6.25 105 radians, which is equal

84
Figure 6.7: Angular resolution of an antenna with aperture D. The red dashed and the black
dashed wavefronts show that isotropic radiators at the top and bottom of the aperture are out
of phase. The maximum of the is on the centerline, in this case all wavefronts originating from
the aperture are in phase.

to 12.9 arc seconds. Do the same for your binoculars and you find 2 arc seconds, which is good
enough to see the Jovian moons. Nobody can do this with the unaided eye. Could your eye or
a small telescope see the lunar landing locations, try this yourself, the answer is negative even
with the best optical telescope on Earth. You would need to fly closer to the moon to be able
to see the lunar landers.
Now back to VLBI. The angular resolution of one VLBI antenna is relatively poor. A typical
observing frequency could be the 21cm wavelength which is equal to the difference in two ground
states of the hydrogen atom. With such wavelengths it is difficult to obtain any form of angular
resolution. The Westerbork array is 2.7 km in length and it can benefit from Earth rotation to
construct a synthetic dish with an aperture of 2.7 km. However, the angular resolution of the
Westerbork array itself is never better than 0.21/2700 or 16 arc seconds which is worse than the
unaided eye.
We can improve the resolution by incorporating a network of radio telescopes, for instance
by combining the observations from Westerbork with Effelsberg in Germany (250 km away) or
Haystack in the United States (more than 5600 km away). The baseline between Westerbork
and Effelsberg will result in an angular resolution of 0.17 arcsecond, and relative to Haystack
one will find 7.7 milli arcseconds which is a phenomenal resolution. Radio astronomy is therefore
able to make higher resolution maps of the skies, but, we are limited to frequencies of natural
radio sources.
Apart from its use in astronomy, radio telescopes are also used for determining Earth ori-
entation parameters and measuring baseline vectors between the VLBI stations. We refer to
figure 6.9 where the differential phase measurement of radio signals received by two radio tele-

85
Figure 6.8: Left: Radiation pattern of a single dipole that is positioned vertically in the center,
the 3D shape would resemble a torus, so it is mostly an isotropic antenna, Right: Radiation
pattern of a seven element dipole antenna for UHF frequencies (in this case 430 MHz). This type
of antenna is called a Yagi antenna, before the invention of cable TVs nearly every household
had such an antenna on a mast on the roof of the house. By aligning a number of dipoles
in phase we can achieve a gain of approximately 15 dB relative to an isotropic antenna, the
standard dipole would only achieve 2.11 dB in theory.

scopes is illustrated. Imagine a radio source somewhere in some galaxy which is seen by telescope
B in figure 6.9, a little moment later in time the signal will be picked up by telescope A. In VLBI
we measure the time difference t by cross-correlation of the received signals from stations A
and B. Both stations are equipped with a hydrogen maser and the correlation of signals involves
mailing storage units with digitized astronomic data to a central processing center, the Jive
facility in the Netherlands is such a center. Since the position of the radio-source (formally we
speak about right ascension and declination) is known we can approximate the orientation of
the baseline between A and B relative to the radio front. If the angle between the baseline A B
and the wavefront is called then the relation between the baseline distance D and t becomes
ct = D sin() so that we are able to determine the length of the baseline A B. In reality the
estimation process is somewhat more complicated because, as said before, Earth orientation
and length of day parameters will are part of VLBI data processing. As of today, VLBI is one
of the main information sources used by the IERS, it is the only observation technique that is
able to measure precession and nutation parameters.

6.3.3 Satellite gravimetry

Starting in 2000 there have been three satellite missions entirely dedicated to the observation of
the Earths gravitational field. Before this date high resolution information about the gravity
field had to come from analyzing tracking data of a variety of satellites, altimeters observing the

86
Figure 6.9: Very low baseline interferometry (VLBI)

87
ocean surface, or terrestrial gravity anomalies measured by so-called gravimeters. The cause of
spatial variations in the gravity field is mostly caused by density contrasts within our planet.
These density anomalies may appear deep within the Earth where they affect the low degree and
orders of a spherical harmonic coefficient set that describes the properties of the Earths gravity
field. Densities closer to the surface generate higher degree and order effects. Finally there is
the rotation of the planet, this causes a permanent deformation of the body and it results in one
term in the gravity field, hereafter we will call it the C 2,0 term, that is substantially larger than
all other terms in the field. To summarize the discussion, the Earths potential V is described
by:
Lmax 
e X ae n+1 
V = C nm cos m + S nm sin m P nm (cos ) (6.8)
ae r
n=0

where e is the gravitation constant of the Earth, ae is the mean equatorial radius, C nm and
S nm are potential coefficients, denotes co-latitude (/2 where is latitude) and longitude
of the observer. Expression (6.8) satisfies the Laplace equation outside all attracting masses,
the normalized associated Legendre functions where already discussed in chapter 3. The role
of equation (6.8) is that the gravitational acceleration felt by a satellite is obtained from the
gradient of V .
Satellite gravimetry experiments focus on the determination of the potential coefficients that
you see in equation (6.8). Preferably we would like to determine the potential coefficients to high
degree and order, however, the reality is that most satellites will not experience accelerations
from terms that exceed degree Lmax = 120. Also, a second problem is that any satellite in orbit
will experience thermospheric drag and solar radiation pressure effects which is a significant lim-
itation. Also, drag and radiation pressure are non-conservative, while gravitational acceleration
is conservative. To solve this problem the solution was proposed to incorporate an accelerometer
in the spacecraft. Any non-conservative acceleration will then be observed by the accelerometer
while gravity will not be observed.
The CHAMP satellite was the first of the three dedicated missions to test this concept, it has
an accelerometer and a GPS receiver, the mission was launched in 2000 and it decayed in 2010,
for details see http://op.gfz-potsdam.de/champ/. The main result of this mission is that
Lmax = 115 was obtained by solving the potential coefficients in equation (6.8) by analyzing the
conservative forces on the spacecraft. One of the main limitations is that CHAMPs accelerom-
eter had a limited accuracy, but also, that only the static gravity field could be observed.
The GOCE satellite was launched in 2009 and it decayed in 2013. The GOCE experiment
was an attempt to improve the resolution of the static gravity field. In this case the satellite is
equipped with a number of accelerometers in a sturdy frame, so that differential measurements
of the Earth gravitational acceleration can be performed. Later we will call this the tidal forcing
of the Earth, the lectures on tides (see chapter 14) will discuss the physics. By measuring the
tidal forces we get the gravity gradients of the terrestrial gravity field and an inversion of these
measurements has resulted in mostly static solutions where Lmax in equation (6.8) is raised to
degree and order 224. With this resolution we can see gravity anomaly structures on the Earth
surface with a resolution of better than 100 km, it allows for instance geophysical interpretation
in areas where terrestrial gravity anomalies could not be observed. For details about the GOCE
mission the reader is referred to the ESA website, cf http://www.esa.int/goce.
The third (still on-going) gravity exploration mission is called GRACE. This mission consists
on two CHAMP like satellites, there are GPS receivers and accelerometers on board of both

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satellites, and the key observation is an inter-satellite range measurement with an accuracy
better than 1 m/s2 . With this set-up it is also possible to measure the gravity gradient of
the Earth, but then at a lower spherical harmonic resolution, approximately up to Lmax = 60.
However, these solutions can be computed within approximately a month so that a time series
of changes in the Earth gravity field can be made. With this information we map the mass
changes in ice sheets and glaciers and determine the mass flux contributing to sea level rise.
Also continental hydrology such as changes in water storage on land can be observed with
GRACE. More details about this mission can be found at the Center of Space Research at the
University of Texas in Austin, cf. http://www.csr.utexas.edu/grace/.

6.4 Exercises
1. The orbit of the Lageos satellite is approximately 6000 km above the Earths surface and
it is entirely determined by satellite laser ranging. The tracking of the CryoSat-2 satellite
is done by laser, but we also have a DORIS receiver on the satellite. The receiver listens
to beacons on the ground that transmit a steady tone at a frequency of 400 MHz and 2000
MHz. Preprocessing is applied to remove all refractive effects from the tracking data: i)
which effects are modelled in a refraction model? ii) what information do you need for i?,
iii) what does the ionospheric correction look like?

2. GRACE is a satellite gravimetry mission that has observed ice mass loss over Greenland,
Antarctica and most other ice surfaces. Explain how precision orbit determination is
relevant for the mission.

3. JASON-3 is an ocean altimeter which observes, like any other altimeter, height change
estimates of the ocean from which we can derive the rate of sea level rise. Reflect in max
15 lines on the question whether the altimeter observed rate of change is compatible with
GRACE results.

4. Eventually GRACE and JASON-3 are going to produce change maps for researchers.
Explain whether both maps will reveal the same spatial details?

5. Explain why VLBI is the only technique capable of determining a precession and nutation
effect, while other techniques can not.

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Chapter 7

Observations, models and

parameters

All information that we acquire from satellites can be called data, this could refer to information
that we retrieve by telemetry from the spacecraft bus collected from the subsystems. The
information could therefore come from star trackers, sun sensors, reaction wheels, the orbit
control system, voltages currents and temperatures. But it could also refer to the status of
processors within subsystems, actuators or ultra stable oscillators. All this information is usually
referred to as housekeeping data, it tells us something about what the satellite is doing and how
it behaves in a space environment.
The payload on a scientific satellite usually consists of instruments such as spectrometers that
operate in various parts of the electromagnetic spectrum, cameras, magnetometers, accelerom-
eters, radar altimeters, laser ranging instruments, or synthetic aperture imaging systems. But
also there is tracking data which is required to eventually determine the position and velocity of
the spacecraft relative to a tracking network on the ground. Also, nowadays many satellites are
equipped with global positioning system receivers, the data is either used by the AOCS (Attitude
Orbit and Control System) of the spacecraft, or it is relayed to the ground. All satellites are
monitored from the ground and commands are generated and verified within the ground control
segment before they are sent to the satellite. The commands that go to the satellite could also
be called data, it consists of orbit and attitude maneuver instructions but also instructions to
control heaters within the spacecraft as well as instructions for onboard instruments. Some well
known organizations that are concerned with telemetry (literal translation: to measure at a
remote location) and control of satellites are the European Space Operations Center (ESOC) in
Darmstadt Germany but also the Goddard Space Flight Center (GSFC) in Greenbelt Maryland,
USA, or the Jet Propulsion Laboratory (JPL) in Pasadena California, USA.
During the lectures on satellite orbit determination we will mostly focus on the retrieval and
processing of satellite tracking data, although we will also consider data of a selected number of
instruments. The first task in the processing of data is usually to replicate the behavior in the
best possible way on the ground. There are various reasons why this is helpful, one of them is
that replication assists in the identification of parameters that qualify the data. Another is that
instruments in a spacecraft often behave differently compared to their pormance on the ground.
So far we have encountered three essential concepts, namely observations, models and param-
eters. Observations are the data that we receive from instruments, systems or sub-systems

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on the spacecraft, models are procedures that depend on a mathematical description, a model
is always implemented as a computer algorithm, and parameters are the control knobs of the
model. By definition a model approximates the reality, and the reality is close to what the
observation data looks like, model parameters can now be optimized so that the model output
resemble the observations. The word data appears everywhere, in daily practice it can either
refer to observations or model parameters. To avoid this ambiguity we will assume that:

y = A(x) + nl y = Ax + l (7.1)

where y is a vector of observation data, and A is a matrix with on each row a observation
equations and x is a vector of parameters that control the degrees of freedom of the non-linear
model A(x). Vector  (regardless whether it belongs to the non-linear or the linear model)
contains the difference between the observations and the model output, and it is often called
a vector of residuals. In this form eq. (7.1) is in principle non-linear which means that the
relation between the observations and the parameters is non-linear. It does not mean that the
observation equations are entirely linear, this concept will be explained later on in these lecture
notes. Typical tasks that we will consider within the scope of SDP are:

Collect and process tracking data from a ground station to a satellite to reconstruct and
predict the trajectory flown by the satellite.

Acquire the attitude data generated by the AOCS, and reconstruct the orientation of the
satellite in space to high precision

Assess the performance of hardware such as the clock oscillator on the satellite, and verify
whether it depends on the temperature

This summary shows that different types of observation data exist, some of the observation
data is directly from satellite instruments while other observation data is already part of a sub-
system that may contain its own control loop. But more importantly, some of the observation
data tells something about the behavior of a dynamic system, while other observation data
doesnt because it tells something about for instance a hardware component. At this point we
introduced another concept, what is a dynamical system, as will be explained in the following.

7.1 Dynamical systems and statistics

Within the scope of SDP all dynamical systems can be described by ordinary differential equa-
tions (ODEs) so that there is a set of first-order ODEs in the form of:

u = F (t, u) + G(t) (7.2)

whereby the state of the system u(t) depends on time t. Furthermore there is an initial state
u(t0 ) = s0 whereby t0 refers to an initial time (or epoch) t0 . The exact shape of the functions
F and G and the state vector u(t) depends on the problem (i.e. scientific or technical case) that
we are investigating.
Dynamical systems have the ability to predict the future state of a system whose initial state
is known at t0 . However, the realism of these predictions depends on at least two issues, namely
whether the state vector u0 is accurately known, and secondly, whether F (t, u) is adequately

91
formulated. Any error in either component will affect the prediction accuracy and cause the
future state-vectors to deviate from reality. In this context it is therefore desirable to be able to
adjust the dynamical system so that the output u may be confronted with observations y as in
eq. (7.1). The exact relation between the state vector u and the model parameters x must be
defined. This is actually a difficult problem and we will take sufficient time in the lectures on
SDP to explain this problem carefully.
But at this moment we can already conclude that we should become familiar with two
different mathematical disciplines, on the one had we deal with ordinary differential equations
as in eq. (7.2) and on the other had we deal with the minimization of a vector of discrepancies
 as in eq. (7.1). The parameter vector x and the observations y will later be related to the
state-vector u. Minimization of the vector  is one of our goals and in this context we will look
into least squares minimization procedures.
We will see that there are several approaches to estimate the parameters x, and this simply
depends on the way we approach the problem. The model parameters could for instance affect
the observations only within a certain time window, or they could affect the observations at all
epochs. We may choose to run the least squares minimization algorithm once (a so-called batch
approach), or we could decide to partition the minimization approach in a sequential approach.
In chapter 8 we will start to introduce concepts required to solve eq. (7.1) in the case where we
are dealing with a linear relation between observations and parameters. We will finish with the
least squares approach for linear problems and also present an approach to deal with non-linear
problems. In section 8.8 we explain that parameter estimation often results in normal equations
whereby the eigenvalue spectrum is positive semi-definite so that some eigenvalues become zero.
In this case we can demonstrate that there is a solution manifold for the parameter estimation
problem. In chapter 10 we go one step further, and we develop a class of parameter estimation
problems that involve a dynamic system approach. At the end of this chapter we present a
sequential method that results in the well known Kalman filter algorithm.

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Chapter 8

Parameter estimation

In eq. (7.1) we show a general approach whereby a model A(x) is able to replicate the observation
vector y where x is a parameter vector controlling the behavior of the model. The difference
vector  will now be our concern, because intuitively we feel that its norm should be minimized.
But before we go into that discussion we present the probability density function of random
variables, next we introduce the expectancy operator including some applications that lead to
the definition of the mean of a data sample, but also concepts such as skewness and kurtosis that
tell us more about the shape of the probability density function. In section 8.4 we introduce a
covariance matrix, in section 8.5 we minimize a cost function that we define for the residual  and
in section 8.6 we extend this approach to non-linear parameter estimation problems. Section 8.7
summarizes properties of the least squares algorithm, advanced topics are in section 8.8 and the
implementation of algorithms can be found in section 8.9.

8.1 Random variables

Vector  in eq. (7.1) is assumed to contain random variables if we assume that a model is
reasonably close to the observations. A random variable is best illustrated with an example,
such as with the outcome of a dice that we roll a number of times. Let X be the number of times
that a 6 appears, and the domain of X is a set of numbers {x = 0, 1, ..., n}. Another example
is that of a point on the Earths surface that is chosen at random. Let X be the latitude and
Y the longitude, and the domain of X and Y is: {X, Y, /2 X /2, 0 Y 2}. These
examples are adapted from pg. 439 in [63].

8.1.1 Probability
Probability is the next topic after random variables, and it is best introduced by the outcome
of an experiment which we define as set S. Now let there be two subsets, hereafter called event
A and B that appear within S as is shown in figure 8.1. The theory of Venn diagrams can now
be used to introduce the concept probability. Let p [0, 1] be a probability function which is
defined such that p(A) results in the numerical probability that event A occurs within S. A
similar situation will occur for event B. By definition p(S) = 1. Two new probabilities will
now appear, one is called p(A + B) and the other is p(AB). The probability p(A + B) is the
chance of an event to occur either within A, or in B and thus by definition p(A + B) occurs
within S. In digital electronics an A + B signal would be the result of a logical OR operation,

93
 Figure 8.1: Probability of two events A and B within the experiment domain S.

while an AB signal would be a logical AND operator. In probability theory the discussion is
whether both events are mutually exclusive in which case the p(AB) outcome would be zero, in
a digital circuit this means that when two random logical signals A and B are fed into an AND
gate that the outcome would always be zero. The opposite would happen if we had mutually
non-exclusive events, in this case the probability that A and B happen at the same time is not
zero. To summarize the situation:

p(A) + p(B) for mutually exclusive events
p(A + B) =
p(A) + p(B) p(AB) for mutually non-exclusive events

8.1.2 Conditional probability

If events A and B are mutually non-exclusive then it is also possible to define a so-called
conditional probability. There are two variations, namely the conditional probability p(A|B) in
case event A depends on event B, and the other way around when B depends on A with the
conditional probability p(B|A). Both conditional events can be related to the Venn diagram in
figure 8.1, the conditional probabilities are as follows:
 
p(A|B) = p(AB)/p(B) A to depend on B
is the conditional probability for event
p(B|A) = p(AB)/p(A) B to depend on A

A special situation occurs when events A and B are independent, in this case p(A|B) = p(A)
and also p(B|A) = p(B). Only in this case we find that p(AB) = p(A)p(B).

94
8.1.3 Bayes theorem
Once conditional probabilities are defined the step towards Bayesian inference is easily made.
Thomas Bayes was an English statistician (1701-1761) whos work resulted in An Essay written
by his friend Richard Price, cf. [47]. Bayes theorem follows from the definition of conditional
probabilities, where the central question is to investigate the probability of event A to depend
on event B which are both subset of the set S. This conditional probability can be derived from
the relation:
p(A)p(B|A)
p(A|B) = (8.1)
p(B)
which is only true when p(B) 6= 0, and as a result we can accept or reject the hypothesis that
event A depends on event B within S. Bayesian inference has numerous implications, an example
is to test the conditional probability that the Truth has a certain probability p(A) and that
your Data has a measured probability p(B), and that you know in advance the conditional
probability that your Data depends on (or says something about) the Truth p(B|A). (In the
16th century this was undoubtedly the most difficult challenge for anyone to accept, namely the
fact that Data and Truth have a probability, and that there are conditional probabilities).
For this example, the probability that the Truth depends on the Data (or better said, that
the Truth is confirmed, or supported by the Data) can be inferred from the Bayes theorem.
p(T ruth)p(Data|T ruth)
p(T ruth|Data) = (8.2)
p(Data)
See also [65], Bayes theorem allows one to infer the reverse question, namely, if we measured
that event B depends on A, and if we know at the same time the likelihood of A and B then
we can apparently also infer the probability that A depends on B. Conceptually a Bayesian
algorithm looks as follows:
p := [p1 . . . pn ]
q := [q1 . . . qn ] where |q = 1|
[p1 q1 . . . pn qn ]
r := Pn (8.3)
i=1 pi qi
where p is a vector of probabilities of type B that depend on type A which itself is distributed
over n channels called Ai , q is the contribution of each channel Ai to A as a whole. The result
of the algorithm is vector r which tells us the probability that if event B happens, that it will
happen on channel Ai .
The only drawback of the approach is the computation of the denominator in eq. (8.3),
because we implicitly assumed that p(B) entirely depends on its visibility in p(B|Ai ) and that
p(B|Ai ) and p(Ai ) are realistically measured or known. In reality Bayesian algorithms need a
training period, it is comparable to the e-mail spam detection problem where one first needs
to see many valid e-mails and e-mails that contain spam. The tokens in the e-mails (usually
words) are then stored in a training dataset, and the conditional probabilities are test to mark
an e-mail as spam or to accept it. For a discussion on this see [54].

Example problem Bayes Theorem

Sometimes it takes advanced reading skills to recognize a Bayesian problem, but a nice example
that I found on wikipedia [2] deals with the probability of failure of products originating from

95
three machines in a production facility:

The situation in the factory is that there are three machines, machine 1 is responsible for
20% of the production, machine 2 is for 30% responsible, and machine 3 contributes
the remaining 50%.

We measured that 5% of the products made on the first machine are defect, 3% fail on the
second, and 1% of fails on the third.

Reversed logic thinking causes the following question. What is the conditional probability
that if something fails, that it will fail on machine i?

And next comes the question that is not necessarily related to the theorem: Are we happy
with this conditional probability, or it is time to put a new machine in place?

For this problem p(A1 ) = 0.2, p(A2 ) = 0.3 and p(A3 ) = 0.5, and also p(B|A1 ) = 0.05, p(B|A2 ) =
0.03 and p(B|A3 ) = 0.01 and we also need p(B) which is the probability of failure for the entire
production facility. Since p(B) was not provided we need to calculate it from what we have:
X
p(B) := p(B|Ai ) p(Ai ) = 0.05 0.2 + 0.03 0.3 + 0.01 0.5 = 0.024 (8.4)
i

The conditional probability that if a failure happens, that it will on a particular machine
follows from the Bayes theorem:

p(A1 |B) = p(B|A1 ).p(A1 )/p(B) = 41.67% (8.5)

p(A2 |B) = p(B|A2 ).p(A2 )/p(B) = 37.50% (8.6)
p(A3 |B) = p(B|A3 ).p(A3 )/p(B) = 20.83% (8.7)

These statistics answer the opposite question that we started with, since our problem that
started with B depending on A where we know the likelihood or contribution of A. This
problem is magically translated into the likelihood that B occurs, and more importantly, that
A depends on B. Bayesian inference is reverse logic thinking which can be very useful in some
circumstances. In the production facility example we could demand that failures have a similar
probability of occurring on each machine, and that it is time to replace the first machine.
The Bayesian algorithm is easily coded in matlab, in the example below vector production
stores the contribution of each machine, vector failure stores conditional probabilities that a
product fails on each machine, vector overallfailure returns the overall probability of a failure
to happen in the factory, and failoneach returns the conditional probabilities of failure by
machine.

% Example of Bayesian statistics in matlab

%
production = [0.2 0.3 0.5]; % known contributions by machine (input)
failure = [0.05 0.03 0.01]; % probability of failure by machine (input)
overallfailure = sum(production.*failure)
failoneach = (production .* failure) / overallfailure

96
As a buyer of products originating from the factory we can now optimize our buying strategy,
because consumer statistics of products informs us about the failures for instance by lot number.
A measure like p(A|B) can then help us to buy those products that have the least chance of
failure, so we dont buy cars that are produced on a Monday morning. Think Bayesian and
youre suddenly smart. Later on in this chapter we will discuss Bayesian parameter estimation
methods where prior information is considered, in this case the theorem is translated in matrices
and vectors.

8.2 Probability density function

Let f (x) be a probability density function (PDF) where x R is associated with event X so
that X [x, x + dx]. The probability p that X [x, x + dx] is now f (x)dx x [x, x + dx]
which is equal to p(x X x + dx). We also Rknow that f (x) 0 and by definition of the area

below the probability density function is one: f (x)dx = 1. As a results, the probability for
an event X to occur on the interval [a, c] is then:
Z c
p(a X c) = f (x)dx (8.8)
a

from which we get that

Z b Z c
p(a X c) = f (x)dx + f (x)dx (8.9)
a b

resulting in:
p(a X c) = p(a X b) + p(b X c) (8.10)
Probability density functions are in fact the normalized histograms that we get from an exper-
iment. The bean machine developed by Sir Francis Galton (1822 1911) cf. [22] is a natural
random generator, and if it is designed properly then all balls will generate a bell-shaped curve
in the bins under the machine that resemble a Gaussian distribution function. The Gaussian
distribution function is a well known PDF that itself depends on the mean and the standard
deviation of x.
1 (x)2
f (x, , ) = e 22 (8.11)
2
associated to this Gaussian PDF is a so-called cumulative probability density function CDF
which is nothing more than the primitive of the PDF:
  
1 x
F (x, , ) = 1 + erf (8.12)
2 2
where erf (x) is a so-called error function. The erf (x) function is implemented in matlab,
and for its mathematical definition I refer the interested reader to literature, cf. [66]. Without
further proof we also mention that Gaussian PDFs are the result of convolving the output of
many other PDFs which by themselves are not necessarily Gaussian, but well-behaved. This
property can be shown by the so-called central limit theorem of which the Galtons bean machine
is a demonstration. Although the mathematical details about the central limit theorem are
interesting the topic goes beyond the scope of what I intended for the class on SPD.

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8.3 Expectation operator
Once the PDF is defined we can continue with the concept of an expected value of X hereafter
called E(X); it is defined as: Z
E(X) = xf (x) dx (8.13)

A property of the expectation operation E(X) is that it easily allows to transform a random
variable X into another random variable via the relation Y = g(X); this results in:
Z
E(g(X)) = E(Y ) = g(x)f (x) dx (8.14)

The corollary of this property is that higher order moments of the expected value may be defined.
As a result we have the so-called k th order moment function of X about 0 defined in the following
way: Z
E[X k ] = k = xk f (x)dx (8.15)

whereby the (weighted) mean or average of X is called 1 . The k th order moment about 1 is
therefore: Z
k = E[(X 1 )k ] = (x 1 )k f (x)dx (8.16)

With the k thorder moment function we can commence to define the variance, the skewness and
the kurtosis of X. The variance is the second-order moment about the mean 1 :
Z
2 2
2 = = E[(X 1 ) ] = (x 1 )2 f (x)dx (8.17)

The skewness 1 is defined as:

E[(X 1 )3 ]
1 = (8.18)
3
and by the old definition of kurtosis 2 is:
4
2 = (8.19)
4
It can be shown that the kurtosis of a Gaussian distribution is 3, and this results in the new
definition 20 = 2 3 so that the new definition of kurtosis should be close to zero when X is
Gaussian. The variance, skewness and kurtosis operators are defined in matlab and you can
call them for a vector with random variables. To investigate the properties of such a vector I
recommend the following analysis:

Plot the histogram (with the hist function) of the data in a random vector, and look at the
difference between the mean and the median functions built-in matlab. If the difference
between both results is large then there are probably outliers in the random vector.

Calculate the second-order moment about the mean with the functions var or std in
matlab. The first question to ask is, does the standard deviation being the square root
of the variance give you what one would expect of the random vector that is analyzed?

98
 Calculate 1 with the skewness function in matlab. If the skewness is far below 0 then
the PDF in the histogram should also look skewed to the left, or when positive then to
the right. Is this also the case in the histogram?

Calculate 2 with the kurtosis function in matlab. If the kurtosis is around 3 then the
distribution is Gaussian, below 3 the distribution is said to be platykurtic, and larger than
three means that the distribution is leptokurtic. Kurtosis will say something about the
peakiness in a distribution. If this is the case then it should be identified in the provided
random vector.

In the following section we will continue with the second order about the mean which is the
variance of a random vector.

8.4 Covariance analysis

The definition of the variance of one random variable X was discussed in section 8.3 and it
resulted in a procedure that can easily be implemented in an algorithm. In case we have more
than one random variable a so-called covariance matrix P will emerge. To demonstrate the
properties of P we select a second random variable Y that is somewhat independent of X so
that the P matrix that is associated with state vector (X, Y )t becomes:

(X E(X))2
 
(X E(X))(Y E(Y ))
P =E (8.20)
(X E(X))(Y E(Y )) (Y E(Y ))2
or:  
XX XY
P = (8.21)
XY Y Y
The elements on the main diagonal of P contain the variances of X and Y respectively, P is
by definition symmetric and the off-diagonal components contain the co-variances between the
variables. At the same time we can now define the correlation coefficient XY between X and
Y:
E[(X E(X))(Y E(Y ))] XY
XY = p = (8.22)
2
E[(X E(X)) ]E[(Y E(Y )) ] 2 X Y
from which we conclude that 1 XY 1. By definition correlation coefficients are symmetric
so that XY = Y X . The covariance matrix of a vector with n random variables takes the
following shape:
11 12 1 2 . . . 1n 1 n
12 1 2 22 . . . 2n 2 n
P = (8.23)

.. .. .. ..
. . . .
1n 1 n 2n 2 n . . . nn

8.4.1 Covariance matrices in more dimensions

With the availability of a covariance matrix of dataset D we can analyse the properties of a
process that is contained in D. In this case D is a datamatrix of m rows by n columns. In
each column vector dj one registers all realizations of random variable Xj and in all rows i one
collects a random vector (X1 . . . Xn ) that is collected at epoch (or event) i. A row vector could

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result from a questionnaire, a test or a measurement that contains n questions or in the case of
a measurement, properties, while the population size (like number of participants or the number
of epochs or events) is m. Let the average for question or property j now be defined as a variable
j derived from the column vector dj in D as j = E(Xj ) where all entries in dj are realizations
of the random variable Xj . Next we calculate the reduced matrix Dr :
 
Dr = d1 1 d2 2 . . . dn n (8.24)

The co-variance matrix associated with (X1 . . . Xn ) becomes:

1
P = D t Dr (8.25)
m1 r
One of the interesting aspects of co-variance matrices is that they contain principle axes, which
are the eigenvectors of P . All eigenvectors are by definition orthonormal and allow one to
construct an ellipsoid. The eigenvectors and eigenvalues of P appear in the following relation:

P uj = j uj j [1, n] (8.26)

so that:
P U = U = U U t = U t U (8.27)
since P is symmetric. For this problem in two dimensions we can plot the columns in the
reduced data matrix in the form of datapoints so that we define an ellipsoidal confidence region.
Figure 8.2 shows D where all green samples fall inside the confidence region, the blue samples
are beyond the ellipsoidal region and have a low probability, whereby the suggestion is raised
that the blue datapoints are anomalies. Once a variance of X relative to its average is known
one can always identify a probability interval within which a majority of the samples are located.
Such an interval follows directly from the CDF of a Gaussian probability distribution function.
For instance, if the standard deviation of X is provided as then the CDF assigns a probability
of 0.68 or 68% to the event whereby samples of X occur in the interval [ , + ]. Confidence
intervals (CI) are usually specified as k intervals, for k = 2 we get CI=95% probability and for
k = 3 we find CI=99%.

Application of confidence intervals

In a manufacturing process, confidence intervals may be used to check the quality of a product.
Coins made by the Royal Dutch mint should have a precise weight and shape; coins forged
by the mint pass a quality control procedure that measures a number of variables of the coin.
The sample to be tested should only be accepted when its measured dimensions qualify certain
production criteria. But this procedure can not prevent that one in every so many coins1 does
not satisfy the production criteria. The region is ellipsoidal because the confidence radius in the
X,Y plane depends on the eigenvalues along both primary axes which are the eigenvectors of P .
Eigenvector analysis works well for 2-dimensional problems, but soon becomes too complicated
when more variables are involved. In 3 dimensions we can still define an ellipsoidal body repre-
senting the confidence region, but in n > 3 dimensions we get hyper-ellipsoidal surfaces which
are difficult to interpret or analyse.
1 100CI
With L 100
times the population size.

100
Figure 8.2: Principle axis of the covariance matrix P and the 3 confidence region. The popula-
tion size n is 5000 for this example, and the confidence region is the 3 relative to the principle
axis.

101
8.4.2 Empirical orthogonal functions
For a general m > 1 by n > 1 with m n dimensioned dataset D there are alternative methods
to investigate the properties of its covariance matrix. One of these methods is to investigate the
so-called empirical orthogonal functions, short EOFs, of the dataset. In the previous section we
introduced D which was reduced to Dr so that the column vectors are centered about 0. We can
subject Dr to a so-called singular value decomposition or svd which comes with the following
properties:
Dr = U V t (8.28)
where both U and V are orthonormal matrices so that U t U = I and V t V = I and where
is a diagonal matrix that contains singular values. The matlab programming environment
has its own implementation of the svd algorithm that you invoke by the command [U, , V ] =
svd(D, 0). The matrix Dr may now be approximated by U 0 V t whereby we zero out some of
the small singular values of . The uncompressed reduced data-matrix is therefore:

1,1
..

.
t
Dr = U n0 ,n0 V
(8.29)

. ..

n,n
where Dr is a m n matrix, U is also a m n matrix and and V are both n n matrices.
When we apply compression on Dr we get Dr0 :

1,1
..
.

0
0
n ,n0 t
Dr = U V (8.30)
0

..
.
0

The first n0 singular values of contain a percentage of the total variance. To understand this
1
property we should look at the relation between Dr and its covariance matrix P = m1 Drt Dr .
t
Substitution of the singular value decomposition of Dr = U V gives:
1 1 1
P = (U V t )t U V t = V U t U V t = V 2 V t (8.31)
m1 m1 m1
1
In other words, if P has the eigenvalues i i [1, n] then i = m1 2i,i so that i,i =
p
(m 1)i . In addition we see that the eigenvectors of P are stored as column vectors in
V . In the approach shown here we discovered that these eigenvalues are associated with random
variables stored in columns, but, the same method of computing covariance matrices may also
be applied over rows. In the latter case we subject the transposed of the reduced data matrix
to a singular values composition. It is up to the reader to demonstrate that U will now contain
the eigenvectors along the row dimension of Dr .
The U matrix will therefore contain column-wise vectors that may be interpreted as eigen-
functions along the row dimension of Dr while V contains column-wise vectors interpreted as

102
eigenfunctions along the column dimension of Dr . The EOFs therefore come in pairs of two
of such eigenvectors that provide an empirical description of Dr . The EOF method is called
empirical because we do not rely on a physical property to obtain the functions, instead, the
EOFs just appear in the U and V matrices after the svd of Dr . EOFs have many applications
in geophysics, they allow one to compress a noisy dataset and to isolate a number of dominating
eigenfunctions of the data. Oftentimes geophysical processes can be recognized in these dom-
inating functions, such as subtle long-term variations in the sea-level measured by a satellite
altimeter, or the prevailing annual wind patterns in a meteorologic dataset.

8.4.3 Transformation of covariance matrices

Let x Vm and y Vn and y = Ax. Without further proof we mention that Pyy = APxx At
where Pxx and Pyy denote the covariance matrices of x and y respectively. The linear trans-
formation implies that both covariance matrices are symmetric. Sometimes covariance matrices
are presented as an expectation Pxx = E[x xt ]; transformation of covariance matrices will be an
essential topic in the following sections.

8.5 Least squares method

Let us now assume that a linear relation exists between an observation data vector y Vm and a
parameter vector x Vn and that we also have a linear model A. In addition we state that there
is a vector of residuals  Vm and that there is a covariance matrix Pyy that represents the noise
in the observations vector. Matrix A is often called a design-matrix or an information-matrix,
and each row of this matrix contains an observation equation, or, is the result of approximating
an observation equation.
The least-squares method was invented by the German mathematician Carl Friedrich Gauss
(1777-1855) who applied the method first to predict the position of dwarf-planet Ceres in the
asteroid-planetoid belt. The least squares method comes in various forms, and we will first show
the simplest case where Pyy = I. The least-squares problem starts with:

y = Ax +  (8.32)

whereby we seek a minimum of the cost function J = t . If we substitute  = y Ax in J then

we find:
J = y t (y Ax) xt At (y Ax) (8.33)
In order to minimize J we seek a vector x that minimizes the second term on the right hand
side of this equation, since the first term cant be minimized when Ax approximates y. If we
exclude the trivial solution that x = 0 then:

xt At (y Ax) = 0 (8.34)

which leads to the so-called normal equations:

At Ax = At y x = (At A)1 At y (8.35)

When the normal equations are solved then x is said to be the unweighted least squares solution
on the system y = Ax + . The matrix At A is called the normal matrix; in this case x is

103
called an un-weighted solution because we did not use information contained in the observation
covariance matrix Pyy to compute the solution. In reality one should use this information in
Pyy because it will redefine the cost function J to be minimized for the so-called weighted least
squares solution in which case:
1
J = t Pyy  (8.36)
The solution of the weighted least squares problem is obtained in a similar way, we seek the
minimum of J and we substitute  = y Ax in J, which gives:
1 1
J = y t Pyy (y Ax) xt At Pyy (y Ax) (8.37)

Also in this case we only need to consider the second term on the right hand side of the equation,
so that the solution for the weighted least squares problem becomes:
1
x = (At Pyy A)1 At Pyy
1
y (8.38)

Note that we may also have found this solution by a reduction operation of the unweighted least
squares problem. The reduction operator is in this case:

y = Pyy
1/2 1/2
y = Pyy [Ax + ] = A x +  (8.39)

This problem may be treated as a unweighted problem because E[ t ] = I. If your computer
implementation of the least squares problem doesnt foresee in the availability of a covariance
matrix of the observations then you should simply reduce your observation data and your infor-
mation matrix as in the above equation.

8.5.1 Parameter covariance matrix

Once you have the least squares solution x, the next problem is to find the covariance matrix of
that solution which is the parameter covariance matrix Pxx . One obtains this matrix by linear
transformation of the observation covariance matrix Pyy . To avoid lengthy matrix algebra we
first assume that we have an un-weighted problem for which the solution was:

x = (At A)1 At y = By (8.40)

so that Pxx = BB t which becomes:

Pxx = (At A)1 At A(At A)1 = (At A)1 (8.41)

The conclusion is therefore that we have to calculate the inverse of the normal matrix to obtain
Pxx . For the weighted least squares problem one first applies the reduction operation, so that
the inverse of the weighted problem becomes:
1
Pxx = (At Pyy A)1 (8.42)

Equations (8.38) and (8.42) implement the general solution of the least squares minimization of
the linear observation model y = Ax +  where Pyy is the observation covariance matrix.

104
8.5.2 Example problem 1
Figure 8.3 shows the result of fitting a function v(t) = a cos t + b sin(2.t) + c.t + d through a (v, t)
point cloud that was observed by an instrument that measures voltages v as a function of time
t. By fitting we mean that we look for an optimal choice of the coefficients {a, b, c, d} so that
v(t) approximates the observations made by an instrument in the best possible way. The best
possible solution is a solution that minimizes the residuals between the observations yi made at
epochs ti relative to the instrument readings v(ti ).
The function v(t) is non-linear with respect to t, but this is not relevant (yet) for the least
squares algorithm because the partial derivatives of v(t) relative to the coefficients {a, b, c, d}
are simply linear. The information matrix A for the least squares problem is therefore:

cos(t1 ) sin(2t1 ) t1 1
A=
.. .. .. .. (8.43)
. . . .
cos(tm ) sin(2tm ) tm 1

the observation vector is:

v(t1 )
y=
..
(8.44)
.
v(tm )
and the parameter vector is x = [a, b, c, d]t If we assume that all observations are equally weighted
with an theoretical variance of one, then eq. (8.35) provides the coefficient values for function
v(t). The resulting function may then be overlaid on the observation data for visual inspection
and the error of the coefficients in v(t) follows from the covariance matrix which is inverse of
the normal equations.

8.6 Non-linear parameter estimation

So far we discussed linear parameter estimation problem, whereby linearity referred to the
content of the design matrix A which is uniquely determined for the problem because the content
does not depend on the parameters x to be estimated. But if this were the case then we would
immediately return to the model concept introduced in eq. (7.1). If there is a non-linear problem
then A can be approximated with the initial guess for x which we will call x0 . We need the initial
guess to be able to linearize the observation equations y(x, t) where t is an independent variable
like time.

y = A(x) +  = A(x0 + x) + 

A 
y = y A(x0 ) x +  (8.45)
x  x0

Next the weighted least squares algorithm should be applied to eq. (8.45). The assumption is
now that the partial derivatives of A with respect to the parameters x should be evaluated at
x0 . If the approximation in eq. (8.45) is adequate then these partial derivatives need only to
be evaluated once. The estimation problem is in that case still linear and one can apply the
algorithms discussed in section 8.5.

105
 Figure 8.3: The linear unweighted least squares function fitting problem

 
A  A 
However, if the partial derivatives differ enough from where x1 = x0 + x0
x x0 x x1
then the algorithm should be repeated with x1 as a starting point. To summarize, the non-linear
weighted least-squares parameter estimation algorithm becomes:

i := 0
repeat
1
xi = xi + (At Pyy A)1 At Pyy
1
(y A(xi ))
i := i + 1; xi := xi
until kxi xi1 k tolerance (8.46)

This algorithm usually converges within a number of iterations, yet the convergence speed can
not be guaranteed. In case of no convergence it should be investigated what is causing the
problem, some well known causes are that the initial state vector in combination with the model
do not describe the observations very well, or, that the normal equations to inverted are nearly
singular because the parameter estimation problem was not well defined. Ill-posed problems
will lead to a manifold of solutions of the least squares algorithm, this topic can be found in
section 8.8.

8.6.1 Example problem 2

Figure 8.4 shows the result of fitting a circle function rc2 = (x xc )2 + (y yc )2 through a (x, y)
point cloud that was measured by a scanner instrument that returns 150 points (x, y) in one
instance. We could be dealing with a production facility that needs to check whether metallic

106
 Figure 8.4: The non-linear unweighted least squares circle fitting problem

dishes on a running belt satisfy the production criteria. The problem is now to determine the
parameters xc , yc and rc so that we find an optimized circle function that fits the scanner
observations in the least squares sense. For this problem we have the linearized observation
equations that contain partial derivatives with respect to the parameters xc , yc and rc . At
iteration k in the non-linear parameter estimation algorithm we have the following observation
equations:
.. ..

. .
xc,k
xi 1 0 (xi xc,k )/rc,k
yi = 0 1 (yi yc,k )/rc,k yc,k +  (8.47)

rc,k
.. ..
.k
.
From these equations we see that the design matrix does contain partial derivatives that depend
on the parameter vector. We defined a circle function with a random number generator, with the
defined function parameters we also generated a random point cloud consisting of 150 points.
Without any prior knowledge the non-linear least squares algorithm was able to retrieve the
circle function parameters in approximately 17 steps. This is an example where the convergence
speed isnt investigated very well, because typically for orbit determination we should be able
to find a converged solution in 3 to 5 steps.

107
8.7 Properties of the least squares algorithm
The conclusion so far is that the inverse of the normal matrix becomes the covariance matrix
of the estimated parameters, see section 8.5.1. The consequence of this property is that the
parameter covariance matrix is sensitive to scaling of the observation variances. Also we can
write the algorithm in three different ways. In the following sub-sections we will discuss these
properties.

8.7.1 Effect of scaling

If one assumes that Pyy = I in equations (8.38) and (8.42) with scaling factor then

1 t
x = At A Ay (8.48)

and
1
Pxx = At A (8.49)
This shows that the estimated state vector is not affected by scaling, but that its covariance is
affected. This property suggest that it is difficult to obtain state vector covariances that are free
of scaling effects, or more generally, we need to determine so that Pxx is in agreement with 
of an observation set. Variance calibration procedures have been suggested by [30] and others.

8.7.2 Penrose-Moore pseudo-inverse

Depending on the number of observations and parameters m and n respectively there are three
implementations of the least squares method, we already demonstrated the first algorithm and
the second algorithm is trivial. Yet the last expression is something new.

1 t
x = At A A y = Ky m > n
x = A1 y m=n (8.50)
t t

1
x = A AA y m<n

To demonstrate the validity of the last expression we consider that:

1 t
1
K = At A A At AK = At AK I = 0 K = At AAt (8.51)

1
which can only be applied when AAt exists so that we should demand that m < n. In the
Kalman filter (that we discuss later on) K is the Kalman gain matrix, and in other literature K
is called the pseudo-inverse of A, or the Penrose-Moore pseudoinverse of A which in literature
is also written as A+ . Regardless of whether m > n or m < n the Penrose-Moore pseudoinverse
satisfies the conditions:

AA+ A = A
A+ AA+ = A+
(AA+ )t = AA+
(A+ A)t = A+ A

In matlab there is a general inversion routine for the system y = Ax where m 6= n. In this case
the inverse of y = Ax is obtained by x = A\y; depending on the dimensions of the A matlab

108
will automatically select one of the three algorithms. The Penrose-Moore inverse A+ also exists
for rank deficient systems, but it will generate a manifold of solutions as will be explained in
chapter 8.8.

8.7.3 Application of singular value decomposition

Singular value decomposition of A = U V t directly results in A+ because:

At A = V 2 U t U V t = V 2 V t

At y = V U t y
V 2 V t x = V U t y V t x = U t y
The last step results in the algorithm:

x = V 1 U t y (8.52)

The Kalman gain matrix is therefore K = V 1 U t . Later on in these lecture notes it will be
shown that 1 is only required for those diagonal elements of that are not zero.

8.8 Advanced topics

In chapter 8 we avoided the problem of linear dependency of column vectors within the design
matrix A. We recall the definition of linear dependency ofPa set of vectors ai Rm , which is
that there exist scalars i 6= 0 and ai 6= 0 that result in m i=1 i ai = 0. If columns of A are
linear dependent then At A will become rank deficient because some of the eigenvalues of At A
will be zero. The result is that the normal equations can not be solved in a straightforward
way. Rank deficient normal matrices occur in practice more easily than expected. Sometimes
the inversion of the normal matrix, At A, will simply fail because of linear dependency resulting
in rank deficiency. Also, many iterations in the non-linear inversion algorithm are usually an
indication that the estimation problem is close to being rank deficient.
It should be obvious that At A is symmetric and that its eigenvalue spectrum is positive
semi-definite. This can be shown with the svd algorithm, in fact, we already did this when we
discussed eq. (8.31) which shows that all eigenvalues are positive. Linear dependency between
the column vectors of A will result in the situation that some eigenvalues of At A become 0. The
number of eigenvalues that become zero is equal to the rank defect of the system. The remarkable
property of least squares problems is that rank deficient problems also have a solution which
comes in the form of a manifold of solutions. To explain this problem better we first need to
speak about the compatibility conditions of systems of equations.

8.8.1 Compatibility conditions of Ax = y

Suppose that we want to solve to n n system Ax = y. We apply an eigenvalue decomposition
on A so that the system becomes U U t x = y or U t x = U t y. The transformations x = U t x

109
and y = U t y yield:

1
..
.
x1

y1


m .. = ..

. . (8.53)
0
xn yn

..
.
0

where m < n. This system shows that there will be no conflict for those equations where
i > 0 i [1, m]. But for i [m + 1, n] there will be an inconsistency in case U t y i 6= 0 . The
compatibility conditions of [32] state that the latter will not occur, in fact, they demand that
U t y i = 0 i [m+1, n]. If xh Vnm is within the null space of A then xh = U s = Es Vnm
where s Wnm is a non-trivial but arbitrary vector. In this case:
 
t 0
U [u1 . . . um |E] = U ty
s
 
0
= U t y y t ui = 0 i [m + 1, n]
s
which demonstrates that xh = Es Vnm is a valid solution. All solutions that occur in the
null space of A are now homogeneous solutions of the system of equations Ax = y of rank m
with x Vn and y Vn . The remaining part of the solution is called xp Vm , this solution is
obtained for the remaining part of the system where i > 0 i [1, m]. The general solution
xg Vn of the problem becomes: xg = xp + Es Vn where, by definition, AE = 0. We can
also say that xg Vn describes the manifold of solutions of the system Ax = y of rank m with
x Vn and y Vn that fulfills the compatibility conditions and that has a rank deficiency of
n m. This is a different situation where one would say that A is singular and that no solution
would exist, in fact, as we have shown, it depends on the right hand side of the system Ax = y
whether we can formulate a manifold of solutions.

8.8.2 Compatibility conditions At Ax = At y

A system of normal equations obtained from a least squares estimation procedure that comes
with a linear dependence in the column space of A will result in a solution manifold. To
demonstrate this property we return to the singular value decomposition implementation of the
normal equations where:
At Ax = At y = r (8.54)
where we apply svd to A = U V t so that:

V T x = U t y (8.55)

We recall that xh = Es where E V so that i = 0 i [m + 1, n].

A[0|E] = U V T [0|E] = 0 (8.56)

110
Figure 8.5: The network solution to synchronize 8 atomic clocks, along each segment the time
difference between two clocks is observed.

Does this satisfy the compatibility conditions? Let us investigate this:

rt xh = y t AEs = 0 (8.57)

so that the answer can be confirmed, rank deficient normal equations will always allow for a
solution manifold because the compatibility conditions are fulfilled. This means that we can
apply a svd on A to investigate its singular values. Equation (8.52) may now be used to obtain
a particular solution xp Vm , and the eigenvectors in V that are in the null space can be used
to formulate xh Vnm .

8.8.3 Example problem 3

Figure 8.5 shows a network of 8 atom-clocks A to H that are connected via a high speed
communications network. Between the clocks one can send synchronization information so that
the time difference along the segment is observed2 . The question is now, how do we adjust the
clocks errors by means of the observed time differences, and, what is the absolute adjusted time
at each clock?
The time difference observations between clock i and j are expressed as ij = Tj Ti
where Ti and Tj denote the clock errors at each atomic clock in the network. The quantity
ij is already corrected for cable and channel delays in the network because round-trip values
are recorded. The epochs at clock i and j are respectively Ti = T + Ti and Tj = T + Tj where
T is the reference time. At each clock we measure the time difference through the network thus
including the delays td along each segment:

ji = Tj Ti + td
ij = Ti Tj + td (8.58)
2
This example exists in real life on the internet where hosts communicate via the so-called ntp (network time
protocol)

111
With this information one can eliminate td and extract the difference of the clock errors which
is what we are interested in.
1
ij = (ji ij ) = Tj Ti (8.59)
2
After this pre-processing step the ij values are related to the clock errors in the network as:

BA 1 1
CB
1 1

DC
1 1 T
A
ED
1 1 T
B
F E 1 1 T
C
GF 1 1 TD

y = Ax = = (8.60)
GA 1 1 TE


HA 1 1 TF


1

HB

1 TG

1

HC 1 TH

HD 1 1
HF 1 1

Any attempt to invert the matrix At A in matlab will now fail because there is a column rank
deficiency of 1 for this problem. In this case we calculate A+ via a singular value decomposition,
whereby A+ = V 1 U t . One will see that 8,8 = 0, for the Penrose-Moore pseudoinverse one
can assume: 1
1,1
..
1 =
.
(8.61)
1

7,7
0
The last column of the V matrix contains now the eigenvector that is in the null space of A.
1
This vector will be like E = 2 2
[1 . . . 1]t so that AE = 0. A particular solution of the problem
is in this case xp = A+ y, and the general solution (or the solution manifold) of the problem is
s
xg = xp + 2 2
[1 . . . 1]t where s is an arbitrary scale factor.
This example shows that the network time adjustment problem has a rank defect of one, and
that we can add an arbitrary constant to all clock epochs. In real-life, one clock is assigned to be
the head-master in the network, and in this way all clocks in the network can be synchronized
to the reference time of this master-clock. Several international networks exist that solve this
problem for the benefit of science and the society in general, it results in the International
Atomic Time (TAI) but also the Global Positioning System time.

8.8.4 Constraint equations

In many parameter estimation problems one prefers to avoid specifying a solution manifold
because it can be a laborious activity. Another reason is also that the singular value spectrum
of the design matrix A changes by iteration step in the non-linear version of the algorithm, or it
can be that the transition of singular values greater than zero to smaller singular values that are
close to zero is not that well defined. The discussion depends on the condition number (At A)

112
which is the ratio between the largest and the smallest eigenvalue of At A. From numerical
analysis theory it is known that approximately k digits are lost in solving the parameters x
when (At A) = 10k . Since eigenvalues of At A are the squares of the singular values i of A it
is from numerical point of view better to go through the pseudoinverse operator A+ so that the
ratio of the largest and the smallest singular value determines the condition number. At JPL
this technique was used for the implementation of Kalman filters on old computer hardware in
deep space explorers that had limited floating point operation capabilities.
A more common approach that avoids to specify a solution manifold is to consider new
information in the form of constraint equations. In this case a-priori information about the
parameters comes in the form of a new set of linear equations c = Bx + c which we now take
together with the observation equations y = Ax + y . The system of constraint equations and
the observation equations are combined into one system:
     
y A y
= x+ (8.62)
c B c

Whereby Pyy and Pcc are available, the cost function to minimize becomes:
1 1
J = ty Pyy y + tc Pcc c (8.63)

whereby we assumed that there no covariance information exists between Pyy and Pcc . At this
point one can simply follow the definition of the normal equations and its solution, so that we
find:
1 t 1
1 1 1
x = At Pyy A + B t Pcc A Pyy y + B t Pcc


B c (8.64)
Eq. (8.64) is in literature known as the Bayesian least squares (BLS) approach. The reason is that
an optimal estimation of x follows from information contained in a model A and observations y
for which we considered constraint information. For normal least squares parameter estimation
we would not assume a-priori information, and hence the relation is made to Bayes theorem.

8.8.5 The Levenberg-Marquardt method

One well known variant of the BLS approach is that covariance information is considered for the
parameter vector x where it is assumed that B = I, c = 0 and where Pcc is provided in advance:
1 1 1 t 1
x = At Pyy


A + Pcc A Pyy y (8.65)

which can even be further simplified by assuming that Pcc = 1 I and an initial guess x0 :
1

1 t 1
x = x0 + At Pyy A + I A Pyy (y Ax0 ) (8.66)

This method is known as the Levenberg-Marquardt algorithm (LMA), Levenberg published the
method in 1944, and Marquardt rediscovered the method in 1963. The recipe for obtaining a
solution of an ill-posed least squares problems with the LMA is to find the smallest possible
so that the condition number of the system becomes acceptable.
In eq. (8.66) we can see that a constant is added to the singular values in the Penrose
Moore pseudoinverse, which will automatically raise all eigenvalues of the normal matrix by .
There is also a better method than the LMA, which would be to raise only those singular values
that are zero, to explain this problem we present all different forms next to one another. Suppose

113
that we have a pre-whitened set of observation equations and that we arrive at y = Ax +  and
that we wish to minimize the norm of the whitened residuals ; the least squares solution is in
this case x = A+ y where A+ for = 0 is:

1
..
.

m
A+ = U
t
V (8.67)

0

..
.
0
In case we would implement the Levenberg-Marquardt method so that the pseudoinverse oper-
ator becomes:
1 +
..
.

+
m + t
A =U V (8.68)


..
.

from which we conclude that the entire spectrum is raised by . This method differs from the
assumption that we only add constraint equations in the null space of A in which case we get:

1
..
.

+
m t
A =U V (8.69)


..
.

which is known as the minimum constrained solution. The benefit of this approach is that we do
not hurt the information contained in the observations by consideration of a-priori constraints.
However, for many problems the reality is that here is no sharp transition between the well-posed
problem where you could avoid any constraints and the ill-posed problem where some singular
values gradually run towards zero. In the latter case one simply attempts to find suitable values
for in the LMA so that the problem appears as numerically stable.
One possibility to implement an approximated minimum constrained solution is a ridge-
regression technique whereby in the LMA is found by seeking the proper balance between xt x
and ty y in J. Ridge regression may be as simple as is to testing different values of and to
plot xt x and ty y for each in an x-y plot. The typical L-shape then appears, and the optimum
is oftentimes found at the corner of the L-shaped curve.

8.9 Implementation of algorithms

There are some standard algorithms to solve least squares problems. Usually we deal with ap-
proaches where apriori knowledge of the observations comes in the form of a diagonal covariance

114
matrix, so that the normal equations can be generated on the fly. Also, several elements in the
design or information matrix A may be zero, so that we can exploit sparse matrix techniques to
invert the equations. For larger parameter estimation techniques blocking methods can be used.
It should also be mentioned that conjugate gradient iteration methods are very effective for
solving large problems if proper pre-conditioners are available. All these techniques fall under
the heading implementation and are discussed hereafter.

8.9.1 Accumulation algorithm, solve on the fly

If the observation equations y = Ax +  are provided with a diagonal covariance matrix Pyy =
then it is trivial that each row and hence every observation can be processed sequentially. Let
yi denote the ith observation in a set of many, and let ai be a sparse row vector of A, and the
normal equations At Pyy Ax = At Pyy1 y can be written as N x = r whereby:

X
Njk := Njk + aij aik /i
jk
X
rj := rj + aij yi /i (8.70)
j

The sums over j and k need to be evaluated for aij 6= 0 and aik 6= 0, and in addition we only need
to store 12 n(n+1) matrix elements which is about half of the elements in the nn normal matrix
N because it is symmetric. After youve processed all observations with this algorithm there is
an equation solver, next one runs again along the information matrix to evaluate the residuals
y Ax. For the inversion algorithm Choleski decomposition as discussed in [46] is popular
because it yields the covariance matrix of the parameters. If there is a Bayesian approach then
an inverse of the Pcc matrix is added to the N matrix and possibly the right hand side vector r is
updated. The accumulation method is popular because intermediate solutions can be computed
while we are processing the observations, at the expense of temporarily setting N and r aside,
hence the procedure is often referred to as solve on the fly.

8.9.2 Sparse matrix solvers

In matlab you have the possibility to store the design matrix A as a sparse structure. The
only consideration is that products like At A need to minimize the amount of stored elements
during the equation solving step. Matlab can automatically do this for you, i.e. select the
best ordering of parameters so that a minimum amount of memory is filled in. Sparse matrix
solvers are useful for many applications for aerospace problems, but their application is limited
to observation equations that are sparse. Nowadays fast memory for matrix storage is not that
much of an issue as it was 30 years ago, and the overhead caused by sparse matrix administration
can make an algorithm unnecessary slow so that full matrix techniques are used. Sparse matrix
techniques become really efficient when they are applied for solving partial differential equations.
In this case band structured sets of equations appear, LU decomposition as discussed in [46] is
often be used, and the reduction of fill-in is guaranteed.

8.9.3 Blocking of parameters

An adequate organization of parameters can help during an estimating algorithm, parameters
may be put together in groups, and in some cases a group of parameters may be eliminated so

115
that the estimation problem remains tractable. An example in satellite orbit determination is
that there are arc parameters and global parameters. The first parameter type in the function
model is related to the set-up of the dynamical model of one arc, and these parameters may
be eliminated by back-substitution (explained hereafter) so that the equation system is reduced
to the set of global parameters. If another arc is computed then the same approach may be
implemented again, until all arcs of a satellite are processed.
In order to implement this technique we assume that the normal matrix can be separated in
four blocks, and that the parameters are partitioned in two sections, namely section x1 and x2 .
In addition we assume that the accumulation algorithm as in eqns.(8.70) already resulted in the
normal equations.
    
N11 N21 x1 r1 N11 x1 + N21 x2 = r1
Nx = r = = (8.71)
N12 N22 x2 r2 N12 x1 + N22 x2 = r2
1
We can multiply the first equation by N12 N11 and to add it to the second so that x1 disappears,
resulting in an equation for x2 , and a similar operation can be performed by multiplying the
1
second equation by N21 N22 and to add it to the first so that x2 disappears. If we assume that
either N11 or N22 can be inverted then we must be able to reduce the system into two separate
equation groups.
1 1
(N11 N21 N22 N12 )x1 = r1 N21 N22 r2 (8.72)
1 1
(N22 N12 N11 N21 )x2 = r2 N12 N11 r1 (8.73)

Suppose that N x = r was a large problem, and that we just processed a batch of observations
where both x1 and x2 appear in the observation equations, but that after this batch of obser-
vations the parameter set contained in x1 will not appear in the observation equations. If this
is the case then we may as well solve eq. (8.72) and continue with eq. (8.73). In that case the
solution for x1 is said to be back-substituted in eqn. (8.73). The above described method may
also be extended over more partitions of N in which case it is referred to as the Helmert-Wolf
blocking method, cf. [12]. Blocking methods can be implemented on distributed computer sys-
tems, and allow one to investigate huge parameter estimation problems. An example is the
parameter estimation problem for highly detailed lunar gravity models where the observation
data is collected by NASAs GRAIL mission, cf. [18].

8.9.4 Iterative equation solvers

Suppose that we have a set of observation equations y = Ax +  where we minimize t , and
where y Vm and x Vn for m n with column rank (A) = m. Earlier in this chapter
we said that x = A+ y which can be a computationally intensive task. The effort to compute
x = A+ y is O(n m2 ) for the accumulation algorithm and O(m3 ) for solving the system of
normal equations, furthermore we did not count the number of operations to define A and y
which can be substantial as well.
The notation O(n) means that the algorithm needs to execute of the order of n operations
to come to an end. The simplest example is an inner product between two vectors, in this
case we need to multiply two numbers at a time and add the result to a sum, this is an O(n)
operation because we simply count all multiplication and add operations in each step as one,
so formally it is n multiplication and add operations, but, we are more interested in the log10

116
of that calculation than the exact number, hence the notation O(n) for the inner product
calculation. An algorithm is well behaved if it can be executed in polynomial time, hence O(nk )
where k > 0. Some algorithms may be optimized from O(n2 ) to O(n log n), sorting algorithms
are a nice example. Sometimes an algorithm may be optimized from O(n3 ) to O(n2 log10 n)
as is the case with two dimensional Fourier transforms or the calculation of grid covariance
matrices from spherical harmonic expressions. Yet there remain a number of algorithms, like
the traveling salesman problem3 , which is of order O(n!). Alternatively, the number of moves
in a game of chess depends on the search depth, recently it was estimated to be at a whopping
O(10123 ) according to [1]. Of course this has triggered the development of efficient algorithms
to minimize the number of search operations. It could be a nice topic to study, but this is not
what we are after in this class.
For very large sets of equations that depend on many parameters the question is whether
we should try to calculate A+ at all, because the problem may be expensive even in polynomial
time. In some cases the exact pseudo-inverse of A is not necessary, so that one we can live
with an approximation of A+ . A simple example of an iterative inversion scheme is to solve the
system of equations y = Ax = (I + L)x with A being a positive definite n n matrix. The
inverse of A looks like (e.g. just develop A1 (L) = 1+L
1
as a Taylor series around L = 0):

A1 = I L + L2 L3 + O(L4 ) (8.74)

so that the solution of the system may be approximated in an iterative approach where we start
with x0 :

x1 := y L x0
x2 := y L x1
..
.
xi := y L xi1 (8.75)

until |xi xi1 | <  so that we converged or until i > threshold in which case the algorithm did
not converge. In reality eq. (8.75) has only a few applications since there are restrictions on the
condition number of A (which is the ratio between the largest and smallest eigenvalue) and on
the eigenvalues of A. There are better methods for iteratively solving systems of equations, one
of them is the so-called conjugate gradient method which locates the minimum of a function
f (x):
1
f (x) = c y.x + x.Ax (8.76)
2
see also [46]. The minimum of f (x) can be found by following the path of the steepest descend
of f (x) along its local gradient f . This gradient is defined as f = Ax y and one can ask
in which direction one should move if we change x by a small increment x. Suppose that we
previously moved in the direction u and that we want to move along v towards the minimum.
If the second move along v does not spoil the motion we had along u then we must move in a
direction perpendicular to u. This means that:

0 = u.(f ) = u.Av (8.77)

3
How many paths exist along which a salesman may travel to visit all his customers? Think about it, and you
will see that n! paths exist for n customers

117
With definition (8.77) one says that u and v are conjugate vectors. The consideration of both
vectors leads to a number of iterative algorithms, such as the conjugate gradient method and
its nephew the preconditioned conjugate gradient method which takes into account approximate
knowledge on the inverse of A. Without further discussion on the details of the CG method we
present here the standard version that solves the system Ax = y. With the theory in cf. [46]
one can show that the CG algorithm takes the following shape:

r := y Ax
p := r
lold := rT r
condition := False
while (!condition)
:= rsold /(pT Ap)
x := x + p (8.78)
r := r Ap
lnew = rT r
condition := lnew < 1010
p := r + (lnew /lold )p
lold := lnew
end

Although algorithm (8.78) looks more difficult than (8.75) it must be said that the CG method
generally converges faster towards the minimum of f (x) in eq. (8.76). In several applications the
CG method is attractive because one only needs to compute A without having to store it. Also
during least squares minimization the CG method is easy to adapt, because one can replace A
by At A and y by At y in eq. (8.78). The drawback of all CG methods is that poorly conditioned
A matrices easily lead to slow convergence so that the benefits of the algorithm are easily lost.
If we know that A is well behaved, for instance because it comes from a differential molecule
applied on a mesh used for discretizing a partial differential equation, then the CG method
might work directly. But otherwise the preconditioned CG algorithm may lose its attractiveness
because one needs to provide a pre-conditioner matrix which is problem specific.

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Chapter 9

Modeling two-dimensional functions

and data with polynomials

This chapter focuses on the problem to approximate functions and data, it is inserted in these
notes as a comment on the least squares method to fit a polynomial to data which is an exercise
in the class on Satellite Orbit Determination.

9.1 Introduction
Let f (x) be a continuous and n-times differentiable function where the domain of x is such that
x [a, b]. If f (x) is not continuous, or if one of its higher-order derivatives is discontinuous, then
we may split f (x) into sub-domains on x which by themselves are continuous and differentiable.
The problem is to approximate f (x) by a series of polynomials pn (x) so that:
N
X 1
f (x) cn pn (x) with x [a, b] (9.1)
n=0

where cn are polynomial coefficients of degree n and where pn (x) is a yet to be defined polynomial.
There are two variations of the problem, the first option is to assume that f (x) is known and
that it is continuous and differentiable, the second option is that there are data points (xk , yk )
with k [0, K 1] in which case xk [a, b] represents an independent variable (such as time)
at which yk is collected.
Fitting means that you find a function that approximates a known function or measurement
data, how the approximation is realized is discussed in section 9.2. We will discuss both ap-
proaches where we start by fitting a polynomial to data points, this is described in section 9.2; an
example where we model Doppler data collected from the Delfi-C3 satellite is demonstrated in
section 9.3. Here we conclude that it is often better to rewrite the Penrose Moore pseudoinverse
into a more stable version whereby we make use of the singular value decomposion algorithm.
Any solution vector that is close to the null space of the involved normal matrix may be ignored,
and this method has superior properties compared to other methods to directly compute the
Penrose-Moore inverse. For details see the article on http://mathworld.wolfram.com.
With measurements one has to take the data points the way they came out of an instrument,
and hence the quality of the fit will directly depend on the gaps in the data. In the nominal

119
situation we expect that measurements are provided every second for instance, but in reality
you will see that some data is missing, the datagaps may then be a concern and this needs to
be investigated.
In section 9.4 we discuss another variant that is by definition not affected by datagaps
because we can select our own support points to model a function. One may select exactly as
many support points as there are coefficients in the polynomial so that we need to invert a K K
matrix. Oftentimes there is no need to invert a matrix because the solution follows directly from
the problem, the Lagrange polynomial function fitting problem is an example that should be
called an interpolation problem because the series of Lagrange polynomial basis functions will
exactly reproduce the data points. An alternative method is based on the orthogonal properties
of Chebyshev polynomials to approximate a known function, the benefit of this method is that
it minimizes the maximum error between the function to approximate and the provided input
function, an example of the last method will be discussed in section 9.5.

9.2 Polynomials to fit data and functions

A straightforward approach is first to define a number of regularly spaced support points xk
[a, b] and to evaluate yk = f (xk ). Let dx = (b a)/K and xk = a + (k + 21 )dx for k [0, K 1]
be an example on how we could chose the support points and let us attempt to minimize  in
the following expression:
N
X 1
yk = cn xnk + k (9.2)
n=0
or alternatively:
1
x20 . . . xN

y0 1 x0 0 c0 0
N 1
y1 1 x1 x21 . . . x1 c1 1
= + (9.3)

.. .. .. .. .. .. ..
. . . . ... . . .
1
yK1 1 xK1 x2K1 . . . xN
K1
cN 1 K1

which we condense to the matrix-vector notation:

y = Hc +  (9.4)

What we can minimize is the L2 norm of the vector . In this case the solution of the problem
follows from least squares minimization where we assume that the variance of y is equal to 2 I.
In that specific case you find for the coefficients: c = H + y where H + is a Penrose Moore inverse
of H so that H + = (H t H)1 H t .
The problem with calculating H + is that the numerical calculation is affected by the scale
of the elements in H but also by the linear dependency between the column vectors that form
H. If there is a linear dependency between any two column vectors in H then the rank of the
H t H matrix will directly become smaller than the number of polynomial coefficients N in the
problem. The way to investigate what is going on is:

Rescale x [a, b] to the interval [1, 1] so that the elements in H, Hij always fulfill the
property that |Hij 1|. For this you introduce a parameter = (x)/(max(x)min(x))
where = (min(x)+max(x))/2 and you use k to replace xk in eq. (9.2) and (9.3). Any kn

120
 will always fulfill the property that |kn | 1 which avoids that the normal matrix (H t H)
is filled with excessive large numbers.

A second remedy is to avoid any calculation of the so-called normal matrix A = H t H

because the condition number of A (it is equal to the ratio of the largest and the smallest
eigenvalue of A which is indicative for the numerical accuracy of an inverse of A) behaves
worse than the singular values stored in the diagonal matrix that appears in the singular
value decomposition H = U V t . After the decomposition matrices U and V become
orthonormal so that U t U = I and V t V = I. The eigenvalues and eigenvectors of the
normal matrix A follow from V , in fact H t H = U V t (U V t )t = V 2 V t which shows
that the singular values stored in are equal to the square roots of the eigenvalues of the
normal matrix A. A singular value decomposition (svd) of H is from numerical point of
view a better starting condition than an eigenvalue decomposition of the normal matrix
A. What does the singular value spectrum look like is one of the first questions to ask
when the H + matrix cant be computed directly.

You can derive H + directly from the svd of the H matrix, it is relatively easy to show
because A1 = V 2 V t and therefore A1 H t becomes V 2 V t V U t which results in
H + = V 1 U t . In other words, for a least squares problem we get:

y = H + c V t c = U t y c = y (9.5)

This system has a number of attractive properties. First of all diagonal elements in
are greater or equal to zero. If they are greater than zero then we can simply invert
the corresponding equations involving c and y . But if you can not invert the relation
because ii < min then we can chose to ignore any corresponding element ci . The reason
is that any ci for ii < min the solutions will appear in or near the null space of A. To
compute a least squares solution you can therefore replace 1 by where ii = 1 ii if
ii min and ii = 0 for all other cases. In this case the least squares solution becomes:
c = V U t y.

9.3 Modeling Doppler data

Let us apply the theory in section 9.2 to a problem where we fit a polynomial function to data
observed with a software defined radio (SDR) listening to the transmitter onboard the Delfi-C3
satellite. The receiver is installed on the roof of the EWI building at the campus of the Delft
University of Technology, from the waterfall plots produced by the receiver we are able to extract
an estimate for the received frequency which contains the Doppler effect of the velocity of the
satellite relative to the receiver. For the provided track the time is specified in seconds. It runs
from 65.5 to 701.5 seconds and the frequency goes from 145.8850095 to 145.8909225 MHz.
Without rescaling the time t [ta , tb ] to [1, 1] I was unable to obtain a polynomial
solution greater than N = 2 in MATLAB. Thus for all results that we summarize in table 9.1
rescaling was applied. The first method in column 2 of table 9.1 shows the standard deviation of
the residuals when you directly compute the Penrose Moore inverse as H + = (H t H)1 H t . The
second method in table 9.1 assumes that H + = V U t where min = 104 . The largest singular
value in the problem is 23.955 and it hardly changes by polynomial degree, the smallest singular
value was 4.42 105 , the ratio of the singular values indicates that the condition number of

121
 N-1 std method 1 std method 2
2 460.057 460.057
3 104.468 104.468
4 89.207 89.207
5 37.201 37.201
6 31.514 31.514
7 19.274 19.274
8 17.769 17.769
9 9.969 9.969
10 10.879 9.449
11 68.090 8.256
12 2622.930 8.199
13 63377.683 7.663
14 582793.602 7.572
15 163889945.529 7.553
16 6927520705.395 7.561

Table 9.1: Standard deviation (std) of the difference between the Doppler track data and a
polynomial consisting of N coefficients, and thus degree N-1 as indicated in column 1. Both
methods are discussed in the text.

the normal matrix can become as large as 2.937 10+11 which means that we could lose up
to 12 digits in any numerical calculation when A is computed directly. The example clearly
shows that the SVD method to compute the Penrose Moore inverse H + is superior to a direct
computation. Figure 9.1 shows the Doppler frequency as observed by the tracking station, and
figure 9.2 shows the residuals of the best fitting polynomial computed with method 2 as shown
in table 9.1. Clearly the residuals show that the measurement noise is not Gaussian distributed,
in the first and last minutes of the dataset there are outliers probably caused by reflections at
low elevations, the center of the residual plot shows a saw tooth pattern which we suspect to be
due to the interpretation of the waterfall plots generated by the SDR. Other effects that could
play a role are frequency variations of the oscillator on the satellite or atmospheric disturbances.

9.4 Fitting continuous and differentiable functions

The theory described in section 9.2 can be applied to any continuous and differentiable function
f (x) on the interval x [a, b]. In fact, it is relatively easy to chose equally distributed supporting
points xk [a, b] and to rescale the H matrix elements between [1, 1]. The algorithm may
be simplified by taking as many support points as there are polynomial coefficients N so that
there is no need to compute a Penrose Moore inverse, instead, a system with N equations and
N unknowns will appear. There is a method to directly solve this problem, that is, given are K
support points (xk , yk ) with xk [a, b] with k [0, K 1] and we seek a polynomial that will
go exactly through the data. The interpolating polynomial is then:
K1
X
L(x) = yk lk (x) (9.6)
k=0

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Figure 9.1: Doppler frequency relative to average frequency measured by the SDR recording a
track from the Delfi-C3 satellite

Figure 9.2: Residual in Hz of the data displayed in figure 9.1 where method 2 was used to
approximate the measured Doppler track data.

123
where a so-called Lagrange polynomial is constructed through the support points:
Y x xm x x0 x xk1 x xk+1 x xM
lk (x) = = ... ... (9.7)
xk xm xk x0 xk xk1 xk xk+1 xk xM
0<m<M
m6=k

and where M = K 1. The only problem with this method is that one can not avoid spurious
oscillations, one possible cause for such oscillations could be poorly chosen support points,
and in particular at the edge of the domain. To circumvent this problem we may chose a
smaller stepsize between the support points at the edge or near a datagap. In this context it
should be mentioned that there are other ways to deal with the oscillations. A remedy is not
to use an arbitrary polynomial like in the Lagrange problem, but an orthogonal polynomial
to approximate f (x). There are numerous orthogonal polynomials but a popular one is the
Chebyshev polynomial function basis. Chebyshev polynomials are defined as:

Tn (x) = cos(n arccos(x)) x [1, 1] and n 0 (9.8)

and they come with a number of properties that interpolating Lagrange polynomials lack. At-
tractive properties of Chebyshev polynomials are: 1) |Tn (x)| 1, 2) Chebyshev polynomials are
orthogonal, 3) there is an equation to directly compute the roots of any Chebyshev polynomial,
4) there are recursive relations to compute the Chebyshev polynomials. For orthogonality:

Z 1 0 : n 6= m
dx
Tn (x)Tm (x) = : n=m=0 (9.9)
1 1 x2 2 : n = m 6= 0

The roots of a Chebyshev polynomial TN (x) follow from the following relation:
 
2k + 1
xk = cos where k [0, N 1] (9.10)
2N

Implementation of the orthogonality relations on the roots for all Ti (xk ) where i < N results in:

NX1 0 : i 6= j
Ti (xk )Tj (xk ) = N : i=j=0 (9.11)
N/2 : i = j 6= 0

k=0

Recursive relations to compute Chebyshev polynomials are:

T0 (x) = 1
T1 (x) = x (9.12)
Tn+1 (x) = 2x Tn (x) Tn1 (x) n 1

Orthogonality relations may be exploited directly to obtain Chebyshev polynomial coefficients

that appear in a series which by itself is meant to approximate an arbitrary continuous and
differentiable function h(x). The task is to estimate the coefficients ci in:
N
X 1
f (g(x)) = h(x) ci Ti (x) where x [1, 1] (9.13)
i=0

124
where y = g(x) is a mapping function to transform the fitting domain of f (y) with y [a, b] to
the required interval x [1, 1]; one possible implementation is g(x) = a + (b a)(x + 1)/2 but
if a different mapping is desired then anything else may be used as well. Let us now integrate
both the left- and the right side of eq. (9.13) in the following way:
Z 1 Z 1 N 1
dx X dx
h(x)Tj (x) = ci Ti (x)Tj (x) (9.14)
1 1 x2 1 i=0 1 x2

By application of the orthogonality properties as in eq. (9.9) we will retain only those polynomials
on the right side of eq. (9.14) where i = j. The right hand side will evaluate to a constant
multiplied times ci . In other words, we have found a way to compute the polynomial coefficients
directly. Drawback is that we need to integrate a continuous and differentiable function h(x)
times all Ti (x) to retrieve the coefficients ci .
For this reason it is more convenient to use the orthogonal relation in eq. (9.11) where we
sum the coefficients over the roots xk of the N th polynomial as outlined in eq. (9.10). Instead
of integrating the left- and right side of eq. (9.13) we can more easily insert the summations (in
fact quadrature relations to replace the integrals), so that we get:

NX1 N
X 1 N
X 1 0 : i 6= j
h(xk )Tj (xk ) = ci Ti (xk )Tj (xk ) = ci N : i=j=0 (9.15)
N/2 : i = j 6= 0

j=0 j=0 i=0

The consequence is that:

N 1
1 X
c0 = h(xk )T0 (xk ) (9.16)
N
j=0
N 1
2 X
ci = h(xk )Ti (xk ) i [1, N 1] (9.17)
N
j=0

where the nodes xk follow from eq. (9.10). In this way we do not need to calculate a Penrose
Moore inverse of the design matrix, we only need to compute the coefficients in eqns. (9.16) and
(9.17) and inspect their behavior as we increase N .

9.5 Example continuous function fit

In table 9.2 we approximate the function h(x) = ex with a Chebyshev series as in eq. (9.13),
next we inspect the coefficients ci for a chosen N to approximate h(x). The magnitude of ci
will indicate the largest deviation in the approximation because |Ti (x)| 1. Table 9.2 shows
that we can approximate ex to within 15 significant digits so that there is no need to take N
beyond 15. The compiler or computer hardware implementation of mathematical functions
usually goes via the evaluation polynomial functions. For this reason Chebyshev coefficients
of known mathematical functions are determined in advance up to a sufficient value of N .
Other applications of Chebyshev coefficients are to compress the results of calculations such
as planetary ephemeris models. For data modeling the Chebyshev function fitting approach
described above one should find a way to first sample the data at nodes xk . This may be a
difficult or problem specific topic that we prefer to keep out of these notes.

125
 i ci i ci
0 1.266065877752009 10 0.000000000550589
1 1.130318207984970 11 0.000000000024979
2 0.271495339534077 12 0.000000000001039
3 0.044336849848664 13 0.000000000000040
4 0.005474240442094 14 0.000000000000002
5 0.000542926311914 15 0.000000000000001
6 0.000044977322954 16 0.000000000000001
7 0.000003198436462 17 0.000000000000001
8 0.000000199212480 18 -0.000000000000000
9 0.000000011036771 19 -0.000000000000001

Table 9.2: Chebyshev coefficients to approximate ex on the domain x [1, 1]

9.6 Exercises
Test your skills:

Rewrite the orthogonal function method to design a procedure where you use Fourier series
to approximate a periodic function. Next investigate how it handles a test function like a
square wave with a duty cycle of 50% between 0 and 2.

Gray function are digital functions that are used for instance in rotary encoders. Gray
functions are also orthogonal. Design your own procedure to transform between the time
domain and the Gray domain.

Demonstrate that empirical orthogonal functions follow from a data matrix subjected to
a singular value decomposition. A data matrix contains measurements or model output of
a defined space that is repeatedly observed. Each observation vector is then stored as a
column vector in the data matrix. Use for instance the RADS database where you select
a repeating track, and show that you can use a limited number of EOFs to describe the
main trends of that track. Find a geophysical interpretation for the results.

What is more efficient: a) to evaluate the square root function in a computer language as
a Newton-Raphson root finder, b) to apply the Chebyshev function fitting procedure?

126
Chapter 10

Dynamical parameter estimation

For all problems considered in chapter 8 we avoided those cases where the dynamics of a problem
is considered. During orbit determination we do estimate the parameters, there is a design matrix
A and there are observations y, but the content of A is often not easily obtained from a linear or a
non-linear set of equations. The reason is that this information must come from a set of ordinary
differential equations which forecast the state u of a dynamic system at future time steps provided
that the initial condition u0 exists at time step t0 . With the Laplace transformation approach
discussed in section 10.1 we can solve a wide variety of ordinary differential equations, identify
the integration constants, and find analytical approximations of u(t). Dynamical parameter
estimation would then be to change the integration constants in an optimal sense so that the
analytical solution fits to the observed data y. Yet oftentimes we resort to numerical integration
of a system of ODEs where the a-priori state of the system and the forcing functions are provided.
In section 10.2 we present the so-called shooting problem whereby we increase to complexity of
the parameter estimation problem for cases that contain dynamical parameters which relate to
ordinary differential equations which can not easily be solved with the Laplace transformation
approach. Section 10.3 works out the details of the shooting problem where numerical techniques
are used, in section 10.3.1 we play two games to show the difference between initial state vector
parameters and other dynamical parameters, section 10.3.2 shows the numerical implementation
of both cases where we distinguish between the state transition matrix and the sensitivity matrix,
in section 10.3.4 we show the regular least squares method where backsubstitution techniques are
commonly used, and in section 10.3.5 we show the benefits of an implementation in a sequential
approach whereby we will present the Kalman filter approach.

10.1 Laplace transformation approach

The dynamical system theory was introduced in section 7.1 where we said that the dynamical
system behaves like u = F (t, u)+G(t) where u is a yet to be defined state-vector. In a dynamical
system where F (t, u) takes the form of the matrix product F u and where F is independent of
u we find as a solution:
Z t1
(t1 t0 ) t (t1 t0 )
u(t1 ) = U e U z(t0 ) + U e e(t1 t0 ) U t G(t)dt (10.1)
t0

with U containing by column the eigenvectors and the eigenvalues of F . But this is a spe-
cific approach that depends on a constant coefficient matrix F . A more common approach in

127
mathematics is to apply the Laplace transformation on an arbitrary function f (t) which directly
results in: Z
F (s) = L {f (t)} (s) = e(st) f (t) dt (10.2)
0
which also has an inverse transform:
Z +iT
1
f (t) = L1 {F (s)} (t) = lim e(st) F (s) ds (10.3)
2i T iT

where = R(s). As you can see, in the time domain we have the independent variable t and in
the Laplace domain this is s, formally both should directly appear as function arguments, but
we use also the short notation F (s) = L {f (t)} for the transform. Laplace transforms have a
number of attractive properties, because, for almost every function we know already the Laplace
transforms and the inverse Laplace transform, furthermore, there are some other properties
which allow you to add, multiply with constants, convolute, differentiate, integrate etc with the
Laplace transformation or its inverse. Also you can work with matrices and vectors in which case
the transforms map onto each element in the matrix or vector. To demonstrate the usefulness
of the Laplace transform we apply it to the left and right hand side of u = F (t, u) + G(t):

L u = L {F (t, u) + G(t)} (10.4)

which becomes:
sL {u} u0 = L {F (t, u)} + L {G(t)} (10.5)
where u0 is the state-vector of the system at t = 0. This becomes:

L {u} = s1 [u0 + L {F (t, u) + G(t)}] (10.6)

from this point onward one should try to move all L {u} terms to the left side and apply the
inverse Laplace transform on the result and on both sides. Sometimes, actually oftentimes, this
means that we need to solve an linear system on equations in the Laplace domain. We can
illustrate what will happen if F (t, u) can be written as a matrix vector expression F u(t) where
F only contains constants. Only in this case we get:

L {u} = [sI F ]1 {u0 + L {G(t)}} (10.7)

The analytical solution of u(t) is now found by the inverse transform applied to the right hand
side: n o
u(t) = L1 [sI F ]1 {u0 + L {G(t)}} (10.8)

Within maple you can easily implement Laplace transforms because they are part of the MTM
package. When used together with the LinearAlgebra package in Maple this provides a pow-
erful tool to handle most ordinary differential equations. In a dynamical parameter estimation
approach we will now be able to identify integration constants for instance in expression (10.6)
or in (10.8) and derive the observation equations with the obtained solutions. The parameters
in the statistical part of the problem are the integration constants, and the resulting analytical
expressions we have found should be differentiated with respect to the integration constants to
obtain the elements for a design matrix.

128
10.1.1 Laplace Transforms demonstrated
Harmonic oscillator
Suppose that we have a system of ODEs u = F (t, u) + G(t) where u = (u, v) is the state-vector:
    
0 n u P cos(t) + Q sin(t)
u = F (t, u) + G(t) = + (10.9)
n 0 v R cos(t) + S sin(t)
At this point we will investigate the analytical solution of this system. What most people will
probably do is to try a number of characteristic solutions and share their results for the integra-
tion constants at conferences or within the company. Hopefully they found all the characteristic
solutions and identified the all integration constants, but, if you want to sure that there are
no other solutions then the known characteristic (company) solutions, then you should apply
Laplace transformations on both sides of the equation.
 
u
L (s) F u = L {G(t)} (s) (10.10)
t
This problem can be reduced to:
      
s n L {u} u0 1 P.s + Q.
= + 2 (10.11)
n s L {v} v0 s + 2 R.s + S.
At this point we invert this equation, and we apply the inverse transform on both sides so that
the general solution becomes (maple did this for me):
  
cos(n.t) sin(n.t) u0
u(t) =
sin(n.t) cos(n.t) v0
 
1 Q. + R.n
+ (cos(.t) cos(n.t))
2
n 2 S. P.n
   
sin(n.t) S. + P.n sin(.t) P. + S.n
+ + 2 (10.12)
n2 2 Q. + R.n n 2 R. Q.n
We can conclude that this expression yields an oscillator with a periodicity of 2/n in a {u, v}
plane. The system is entirely determined by the initial choice of the integration constants u0
and v0 , this part determines the homogenous solution of the system. But, there is also a forced
response when the integration constants {P, Q, R, S} play a role, and they control the remaining
terms in eq. (10.12). The first thing to conclude is that the solution behaves linear with respect to
the forcing terms in the frequency domain, what we mean is that an forced motion at frequency
results in a oscillation at the same frequency. The second thing we note for this problem is that
all forced motions couple back into the natural frequency n, and the third thing we see is that
a natural amplification will occur when n approaches . In that case the oscillator is externally
perturbed by a signal near its natural frequency, and the oscillator will start to resonate. A
fourth thing to note is that eq. (10.12) becomes singular when the system is forced at a constant
signal, this occurs when is zero, because of the constants P and R.
Resonant solutions occur when n = and this requires a separate approach because of the
singularity in eq. (10.12), but this is easy because we simple plug this information into the
Laplace domain, and we let maple do the work:
    
0 n u P cos(n.t) + Q sin(n.t)
u = F (t, u) + G(t) = + (10.13)
n 0 v R cos(n.t) + S sin(n.t)

129
It becomes in the Laplace domain:
   1    
L {u} s n u0 1 P.s + Q.n
= + 2 (10.14)
L {v} n s v0 s + n2 R.s + S.n

so that we find via the inverse Laplace transformation in the time domain:
   
1 P.t S.t + 2 u0 1 Qt + Rt + 2 v0 + nP + nS
u(t) = cos(n.t) + sin(n.t) (10.15)
2 Q.t + R.t + 2 v0 2 P t + St 2 u0 nQ + nR

This solution shows an oscillation at the natural frequency n, but, it also shows a continuous
amplification because there are terms that are linear dependent with respect to time. For this
problem one should directly ask what will happen on the long run, because at some point one
may expect for a physical system (we discussed here mathematics) that the natural frequency n
can not be sustained indefinitely. In other words, the ODE put forward to model our dynamics
is maybe not as realistic as we want it to be.

The gun bullet problem

For this problem we intend to model the motion of a bullet; we deal with a local gravity g and
the motion of the bullet in an orbital plane. The ordinary differential equations for this problem
are actually second order, and need to be reduced to a system of first order ODEs:
      
x 0 A v x
=
y g 2M v y
p
where the x-axis is horizontal and the y-axis goes vertical and where v = x2 + y 2 . In this case
we can only apply the Laplace transformation approach when the ballistic drag is zero, or when
the equations of motion can be linearized as a constant drag problem. For the non-balistic drag
problem you find:

D(x)(0) x(0)
L {x = 0} L {x} = +
s2 s
D(y)(0) y(0) g
L {y = g} L {y} = + 3 (10.16)
s2 s s
and when you apply the inverse Laplace transformation on both sides you obtain the well known
equations:

x(t) = D(x)(0)t + x(0)

1 2
y(t) = D(y)(0)t + y(0) gt (10.17)
2
As soon as ballistic drag is part of the problem there are no easy analytical solutions, in fact, for
this case we only know the homogeneous solution. The one dimensional drag problem does have
a solution, but it is not obtained with Laplace transforms. In order to obtain particular solutions
one should use other methods, and the most often used method is numerical integration of the
variational equations which we define later on in this chapter.

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10.1.2 Define and optimize the parameters
No matter what youve found at this point, a system of ODEs always comes with one homoge-
neous solution and several particular solutions. For the harmonic oscillator we saw that some
particular solutions may exhibit a resonant response towards forcing at the natural frequency
n. But the most important thing of all is that we found integration constants that determine
the outcome of the problem. Also, there may be terms in the solution of the ODE from which
we may suspect that they are not accurate known. In example problem (10.9) we modelled an
oscillator in the (u, v) plane and we found that eqns.(10.12) and (10.15) describe the general
solution of the problem. The integration constants in these solutions are either homogeneous
or particular, but, they allow one to forecast the motion of the oscillator in the (u, v) plane.
The number of parameters depends on the number of forcing terms, for each frequency we
get 4 extra parameters, and we always have the initial state vector (u0 , v0 ). For the gun-bullet
problem we could only approximate the solution and identify x0 , y0 and x0 , y0 as integration
constants for a parabolic orbit that comes with a constant gravity approach. For least squares
parameter estimation the (u, v) plane oscillator would be embarrassingly linear because the par-
tial derivatives between u(t) and u0 are sin and cosine expressions that (only) depend on time.
This is also true for all other integration constants {P, Q, R, S} in both the non-resonant and
the resonant particular solutions. However, as soon as we insert n or in our shopping list of
parameters to be estimated then all that beauty disappears. In that case we should start with
a Bayesian parameter estimation approach and allow that a-priori information for the param-
eters exists. For the gun-bullet problem we always find linear parameter estimation problem,
because the outcome of the motion linearly scales to the integration constants. All together it
is always useful to understand the analytical behavior of a system and to identify the parame-
ters that need to be adjusted. The reality for orbit determination is that we also stop here, in
the real world problems are solved with numerical equation solvers, and the partial derivatives
required in observation equations should be provided in another way than with the knowledge
of analytical solutions.

10.2 Shooting problems

Figure 10.1 contains an example of a so-called shooting problem, it appears in the game of dart
where the task is to find an initial position and a velocity of the dart so that it will hit the
bullseye on a dartboard. The equations of motion of the dart will obey gravity, and all other
forces such as drag plus maybe some wind effects. During the game a player will throw a dart
and observe where the dart hits the board. Let this be the position (xd , yd ) away from the
bullseye at (xb , yb ), the observation is the difference (xd xb , yd yb ). The observed difference
causes the player to adjust the initial angle and velocity, and maybe also the position from which
the dart is launched. In reality the game is slightly more complicated, because the optimum is
not only to hit the bullseye, but also to hit other positions on the board that yield a certain
amount of points during the game, but for simplicity we only look at the problem of hitting the
bullseye.
The method of the dart player is an implementation of the shooting problem whereby he
or she estimates the initial conditions required for the equations of motion of the dart. The
problem is essentially not much different from what we have in precision orbit determination
(POD), because the dart board in POD is represented by satellite tracking data. A least squares

131
 Figure 10.1: Shooting problem in the game of dart

minimum is sought for the initial conditions of the satellite whereby we hope that the calculated
path of the satellite matches the tracking observations in the best possible way.
The mathematical implementation of the shooting problem contains two essential steps. One
is the ability to implement the forward calculation whereby we insert knowledge of the initial
statevector into the problem, the other is that we can relate observations on the dartboard (and
hence at future epochs) relative to the initial epoch. Let the state u(t) in the game of dart
be vector with 6 elements, being three velocity components and three ordinates to describe the
position. Let x be a horizontal component, y goes away from the player and z is vertical upward,
and the corresponding velocities are u, v and w. In that case the equations of motion for a local
coordinate system are:
x u
y v

z w
u =

= 0
(10.18)
u
v 0
w g
which indeed takes the shape of a system of first order ODEs u = F (t, u) + G(t). (The gun
bullet problem is a simplified version of the equations of motion of a satellite, where the inertial
acceleration of the vehicle is nothing more than u = U where U = /r for the Kepler
problem, or where U is expanded in spherical harmonics in case a full gravity model is taken
into account.) The first question is, how do you integrate this system forward in time? The
second question is, how do we get the system of ODE in a shape that a small disturbance at
epoch u(t0 ) = u0 is translated into u(t), thus:
u(t) = (t, t0 )u(t0 ) (10.19)
where (t, t0 ) is called the transition matrix of the problem. The first question asks us to discuss

132
numerical methods for the integration of ordinary differential equations , and the second problem
calls for the integration of so-called variational equations, these problems will be discussed in
sub-section 10.2.1 and 10.2.2.

10.2.1 Numerical integration methods

Without having the intention to be too specific on this topic we present here two methods which
have proven their merits in dealing with solving ordinary differential equations. There are two
type of algorithms, namely single-step methods and multi-step methods for solving a system of
ODEs.

Runge-Kutta method
The RK method(s) deal with the system of ODEs u = F (t, u) for which the initial condition is
u(t). The algorithm returns u(t + h) where h is the step-size of the integrator, without further
evidence we present:
h

u(t + h) = u(t) + k 0 + 2k 1 + 2k 2 + k 3
6
k0 = F (t, u(t))
h h
k1 = F (t + , u(t) + k 0 )
2 2
h h
k2 = F (t + , u(t) + k 1 )
2 2
k3 = F (t + h, u(t) + hk 2 ) (10.20)

This implementation is called the fourth-order Runge-Kutta method because the local truncation
error is O(h5 ). Also, it is a single-step method because it executes four function evaluations to go
forward from epoch t to t + h. Stability of the method can be investigated by consideration of a
suitable test solution, but in practical applications there usually exist (approximated) analytical
solutions for ODEs which may be taken as a reference. Higher-order Runge-Kutta methods
also exist, but they rapidly lose their attractiveness since they require more than 4 function
evaluations while multistep integrators require just two function evaluations to advance from t
to t + h. For precision orbit determination it is however desirable to reduce the local truncation
error to O(h12 ) or even higher depending on architecture of the used hardware. For this reason
we present in the following sub-section a method that can easily decrease the local truncation
error without becoming increasingly expensive on the number of functions evaluations at mid-
points.

Adams-Moulton-Bashforth method
A well known method in precision orbit determination software is the Adams Moulton Bashforth
method. This method needs just two function evaluations to advance a state vector u from epoch
t to t + h and has a local truncation error of O(hn+1 ). This can be achieved by maintaining
a record of earlier function evaluations which is different compared to the single-step method
which performed additional function evaluation between t and t+h. Both Adams-Bashforth, and
Adams-Moulton will not consider function evaluations between t and t + h. The AMB consists

133
of a prediction step in which a future state vector u(t + h) is calculated, and a correction step
whereby use is made of an extra function evaluation at the predicted point u(t + h). The
predictor method is called the Adams-Bashforth method whereby:
n
X
u(t + h) = u(t) + h m F (t + (1 m).h, u(t)) (10.21)
m=1

which shows that each term on the right hand side must be known, in addition the step size
h is fixed for all previous function evaluations. The second part of the algorithm is called the
Adams-Moulton method, which benefits from the knowledge of the predicted state-vector at
epoch t + h. The Adams-Moulton method closely resembles eqn. (10.21) but it includes an extra
step:
Xn
u(t + h) = u(t) + h m F (t + (1 m).h, u(t)) (10.22)
m=0

After prediction and correction a recycling operation occurs where u(t) is replaced by u(t + h)
and also F (ti.h, u(ti.h)) is replaced by F (t(i+1)h, u(t(i+1).h)) for all i [0, n1]. The
efficiency of the AMB method is therefore equal to 2 function evaluations. An one-time effort is
to determine the predictor and the correction coefficients in eqns. (10.21) and (10.22). A separate
problem is the determination of the predictor and corrector coefficients, this is explained in the
following two sections.

AMB predictor coefficients

In order to determine the predictor coefficients m we consider the polynomial:

f (t) = a0 + a1 t + a2 t2 + . . . + an tn (10.23)

of which the first derivative is:

f 0 (t) = a1 + 2a2 t + . . . + nan tn1 (10.24)

An evaluation of f (t) at t = 0 and t = 1 gives:

f (0) = a0 (10.25)
f (1) = a0 + a1 + . . . + an (10.26)

so that f (1) = f (0) + df where df = a1 + . . . + an . Let us now try to determine df for the case
where a1 to an follow from a linear combination of f 0 (0), f 0 (1) etcetera. It is relatively easy
to show that this results in the following system of equations:

f 0 (0)

1 0 0 ... 0 a1
f 0 (1) 1 2 3 . . . n(1)n1 a2


f 0 (2) 1 4 12 . . . n(2) n1
= a3 (10.27)


.. .. .. .. . . .
. .
.
. . . . . . .
f 0 (m) 1 2m 3m2 . . . n(m)n1 an

134
 0
so that f = M a which can be solved by a matrix inversion. Once a is known then:

1 1
t .. 0 t 1 ..
df = a . . df = (f )(M ) . . (10.28)
1 1

at this point the coefficients m required in the corrector method follow from:

1
t 1 ..
= (M ) . . (10.29)
1

In addition, we do not need to worry about the step-size h since it is just a scaling of the same
problem, e.g. f (1) = f (0) + df is the same as u(t + h) = u(t) + h.df . In reality we only need
to calculate the predictor coefficients once and this can easily be coded for instance in maple,
which can directly generate the high-level programming code for you. Some examples are:
 
3 1
u(t + h) = u(t) + h F (t, u) F (t h, u) (10.30)
2 2

and
u(t + h) = u(t) + h {1.916 F (t, u) 1.3 F (t h, u) + 0.416 F (t 2h, u)} (10.31)
where 0.666 denotes 23 .

AMB corrector coefficients

The Adams-Bashforth corrector coefficients follow in a similar way as for the predictor coeffi-
cients. In this case the polynomial is extended so that it also contains the term f 0 (1).

f 0 (1)

1 2 3 n a0
f 0 (0) 1 0 0 0 a1

0

f (1) 1 2 3 n(1) n1 a2
f 0 (2) = 1 4 (10.32)

12 n(2) n1 a
3
.. . .. .. .. .. .
.. ..

. . . . .
f 0 (m) 1 2m 3m2 n(m)n1 an

whereafter:
1
= (M t )1 . ... (10.33)

1
You can verify yourself that:

u(t + h) = u(h) + h {0.416 F (t + h, u) + 0.6 F (t, u) 0.083 F (t h, u)} (10.34)

135
Some remarks on the use of the AMB method
For the integration of the equations of motion of a satellite, and probably also for many other
problems that require one to solve ODEs, one should remind that:

The AMB method requires F (t, u) evaluations at time steps t i.h that occur prior to
t0 and one step after t0 . The consequence is that we first go through a learning period
where earlier F (t i.h, u) values are evaluated by an initialization method. One may use
a fourth-order Runge-Kutta method with a smaller step size so that the integrator error
is sufficiently minimized. Alternatively the RK4 method can run forward and backward
on the AMB initialization points until one is satisfied with the results, or use can be
made of higher than fourth-order Runge-Kutta methods possibly with smaller stepsizes to
initialize. There is no exact solution for this problem, everyone implements his own flavor.

In reality the choice of n in the AMB method depends on the computer hardware, for low
eccentricity orbits one can use an order 11/12 predictor corrector method with a step-size
h as large as 1/100 of an orbital period. An increase of n, or a decrease of h, does not
automatically result in a better orbit. But in case of doubt one can always attempt to
integrate forward and backward to see whether one arrives at the same initial condition.

More than one correction step may be executed in some software, this occurs when there
is a substantial difference between the predicted and the corrected state-vector at epoch
t + h.

For high eccentricity orbits it can happen that the AMB integrator requires a re-initialization
to adjust the step-size h depending on whether arrives at the peri-apsis or the apo-apsis.
Oftentimes the decision to restart the integrator is easier than to focus on variable step
size algorithms.

The Kepler orbit has an analytical solution, and it is easy to verify the quality of the
integrator with the help of the analytical solution. One example has been calculated for
a0 = 7 106 m, e0 = 0001, I0 = 108 , 0 = 0 , 0 = 0 and f0 = 0 , using a step size
of h = 60 seconds. This example shows that the AMB method results in small numerical
errors in the semi-major axis, we find a values like 2.75, 0.53 and 0.02 cm for an order
10/11, 11/12 and 12/13 AMB method. These numbers relate to an arc length of one year in
which case it seems that we are losing (or gaining) energy whereby the disturbances build
up in the flight direction of the satellite. In general, numerical integration errors are small
so that they can be neglected when compared with other modeling effects in precision orbit
determination. The strength of orbit determination really comes from satellite tracking
data in combination with improvement of dynamic model parameters.

The length of the orbital arc to integrate ahead in time is usually confined to either a
week or a month depending on the problem. This is usually not related to the quality of
the orbit integration method, but is rather limited to the reality of the force models that
describe the dynamics of the satellite motion.

And with this conclusion we go to the next section, which deals with the generation of the
partial derivatives in the state transition matrix (t, t0 ) in eq.(10.19) which we need to set-up
a dynamical parameter estimation problem.

136
10.2.2 Variational equations
Now that we have a numerical integrator we can generate trajectories ahead in time. The
dynamical model is u(t) = F (t, u) + G(t) where we start with an initial guess u(t0 ). The
integrator algorithm will now produce a list of statevectors u(t + i.h) where h is a suitable
chosen step-size. We will look into the problem where there is some perturbation u(t0 ) = u0
applied to the initial state, and ask ourself how it results u(t + i.h) for i > 0 and h > 0. One
could simply implement this test in the integrator, apply some perturbation at t0 and see what
happens, but, there is also a formal way and this is where the so-called variational equations
come in view.
The variational equations for this problem are obtained by introducing first a dynamical or
a control parameter k where k [0, K 1] which may appear anywhere within the definition of
F (t, u) or G(t). In the dart problem we could select k to be the local gravity acceleration g but
we may also select one of the components of u. As a result the variational equations become:

u F (t, u) G(t)
= + (10.35)
k k k
which generates 6 first order differential equations for each dynamical parameter; for the 3D
dart problem we end up at 42 differential equations if all elements in u are treated as dynamical
parameters, if gravity is also taken into account then 48 differential equations need to be han-
dled by the numerical integrator, whereby the first 6 elements in the state-vector concern the
equations of motion.

10.3 Parameter estimation

In section 10.1 we presented a method where the dynamical parameter estimation problem was
analytically formulated via (for instance) the Laplace transform approach. The other method
is to use a numerical technique and to formulate the problem as a shooting problem. In the
latter case the partial derivatives obtained on the left hand side of eq (10.35) are obtained by
numerical integration. Regardless of the technique that is used (analytical or numerical) we
obtain a solution for equation (10.19), so that perturbations at the initial state u(t0 ) can
be propagated to any future state vector u(t). Yet the problem extends further than this,
because the variational equations as in eq. (10.35) may also be formulated for terms k that
are not part of the state-vector u in the dynamical problem (7.2). In both cases we are dealing
with dynamical parameters, but, there is a difference and this is best explained by looking at
examples such as the variational problem for the games of dart and curling.

10.3.1 The difference between dart and curling

Figure 10.2 and 10.3 show the basics of both games which have in common that equations of
motion apply to either the dart or the puck which is called the rock in curling. In essence we
can describe the motions as solutions of ordinary differential equations, in dart the dimension
of the state vector is 6 while in curling it is 4. The objective of both games is to reach a target,
but the way the objective is reached varies by game.
In dart the only degree of freedom for the player is to modify the initial state vector, the
gamer may pick any position he prefers as long as he stays behind a line, the direction of the

137
 Figure 10.2: Basics of dart.

Figure 10.3: Basics of the curling.

138
dart and its velocity may be changed and the objective is to hit a target on the board. The
game strategy dictates which segment should be hit on the board. During curling there is a
team and someone throws the stone. It resembles jeux the boules played in France and the
objective is also to reach a certain location on the ice court. The fun part of curling it that, once
the stone is thrown, there are team players with brooms to influence the motion of the stone
while it slides over the ice court. This is the only sport that I know of where this is allowed or
done. The relation to the mathematical shooting problem is that the outcome of the variational
problem u(t) not only depends on the perturbation at the initial position u(t0 ) but, that it
also depends on the influence of other control parameters where k in equation (10.35) differs
from u0
During curling we are interested in the outcome of u(t) that depends on the friction pa-
rameter that is part of the equation of motion of the rock. The horizontal friction may be
approximated by a constant times m.g with m being the mass of the rock and g gravity, also
corresponds to k in eq. (10.35). In this case the variational problem should be defined in a
different way, because the effect of variations in on the outcome of equation (10.35) can only
take effect for t > t0 . This situation corresponds to curling where u(t) is caused by changing
while the stone is in motion. The conclusion is therefore that we should be careful in defining
how eq.(10.35) is treated when we solve the problem, either analytically or numerically. Let us
look at two example problems where the analytical solution is known:

When we discussed the solution of the harmonic oscillator we found solutions as in eqns.
(10.12) and (10.15). Both equations show that u(t) does depend on u(t0 ) but also, that
the outcome is determined by integration constants P through S that were defined for
these problems. Parameters P through S could in this case attain a known constant value,
but, what the variational problem entails is how perturbations P through S affect a
perturbation at u(t). These perturbations do not result in an initial effect at u0 , this
is a necessary condition for the initial value problem.

A second example in equation (10.17) shows the same behavior, any perturbation is g can
only affect the left hand side in case t > 0.

10.3.2 Numerical methods applied to the variational equations

We consider the system of ordinary differential equations of the variational problem:

u(t) = F (t, u(t)) + G(t)
w(t) = A(t, w(t)) + B(t) = u F (t, u(t)) G(t) (10.36)
= + k [0, K 1]
k k k

where w(t) has the dimension 6 + 6K. Within the state vector one finds that u(t) contains
the inertial position and velocity components of the satellite. In addition to u(t) one finds
the partial derivatives u/k , where k . Vector contains the control (or dynamical)
parameters that are defined for the variational problem; there are no specific rules as long as the
partial derivatives exist in (10.36) so that integration is possible. Examples of control parameters
are:
I
Any element within u(t), the corresponding control parameters are called

139
 Any element within that is part of F (t, u(t)) and G(t). Typical examples are a scaling
parameter Cd in a thermospheric drag model or the Cr parameter in the solar radiation
pressure forcing model, or, the gravitational constant in the formulation of gravitational
f
acceleration. In the following these parameters are denoted by .
The solution of eqns.(10.36) therefore consists of two types of partial derivatives which relate
I
the changes (or perturbations) in u(t) to either changes in = u(t0 ) at the initial epoch
f
t0 or a perturbation in one of the forcing models. The expressions are:
f
u(t) = (t, t0 )u(t0 ) + S(t, t0 ) (10.37)
where we notice that the elements of the transition matrix (t, t0 ) and the sensitivity matrix
S(t, t0 ) are obtained by numerical integration of eqns. (10.36). The consequence of the above is
that the initial state vector w(t0 ) to use in the numerical integration method should be chosen
with care. By definition (t0 , t0 ) = I in eqn. (10.37) and moreover S(t, t0 ) = 0 at t = t0 . The
f
interpretation is that perturbations in u(t0 ) map to u0 and that any perturbation can
not result in a u0 6= 0, in fact, for the latter part one needs often a sufficient integration time
until |u(t)| > 0 for t > t0 .

10.3.3 Partial derivatives with respect to the control parameters

After integration of eq. (10.36) we obtain partial derivatives for the state-vector at epoch t
relative to the control parameters. This is usually not what we need in precision orbit determi-
nation, because we do not have information on u(t) at t > t0 , instead we have observations to
the satellite, or from a tracking system on the satellite such as GPS or DORIS. In the following
we consider two hypothetical observation types that yield information about the range, and the
range-rate to the satellite which are both obtained from an observer at the Earth surface. We
seek the partial derivatives of those observations with respect to the control parameters in the
problem, this information is always required in observation equations.

Range observations
Range observations can be obtained in various ways, one possibility is to measure the round trip
time of laser light from a ground based observer to a satellite equipped with a laser retroreflector,
another possibility is to use radio techniques. If we ignore offsets, refraction, relativity and the
light time effect then range observations can be obtained by:
r = |xr xs | = ((xr , xr ) 2(xs , xr ) + (xs , xs ))1/2
where xr and xs denote the position of the observer and the satellite in an Earth-fixed geocentric
coordinate system. Vector xs requires a transformation into an inertial coordinate frame because
we solve the equations of motion in such a frame. For this reason we introduce the transformation
xs = R xi so that the range observation becomes:
r = |xr R xi | = ((xr , xr ) 2(xr , Rxi ) + (Rxi , Rxi ))1/2
which allows us to differentiate with respect to a dynamical parameter k that is part of the
variational problem of the satellite:
 
r 1 xi
= Rxi xr , R
k r k

140
In principle we need these partial derivatives in range observation equations. The term dxr /dk
comes forward from the solution of the variational equations as we have discussed before.

Doppler observations
Doppler observations are nothing more than the range-rate from a terrestrial observer to the
satellite (or the other way around). These observations are obtained by measuring the frequency
change of an oscillator in the satellite relative to an oscillator on the ground. When we ignore
frequency offsets, refraction, relativity and the light time effect then the observation is nothing
more than a projected velocity along the line of sight. The projected velocity can be written as:

v = (v s v r , esr )

where
(xs xr )
esr =
|xs xr |
denotes a directional unit vector with its origin at observer (r) pointing toward the satellite (s);
the partial derivatives are as follows:
v  
= Rv i + Rxi v r , esr
k k
If the vector esr does not change too much then we may assume that:
 
v v i xi
= R + R , esr
k k k

10.3.4 Batch parameter estimation

Suppose that range or range rate observations to a satellite are available and that the problem is
to adjust the dynamic parameters in a variational problem. In addition it should be mentioned
that there may also be non-dynamical parameters in the problem, such as coordinates of the
involved ground stations, earth rotation parameters, and instrumental parameters that correct
the offsets that come with realizing the observations. In principle any of the techniques discussed
in chapter 8 may be used to obtain an estimation for the dynamical and the non-dynamical
parameters that are defined for the problem as long as a sufficient amount of observational data
is collected.

Non-linearity of the problem

An estimation problem for precision orbit determination is by definition non-linear and use is
made of the techniques discussed in section 8.6. The main reason for non-linearity is that the
estimation process starts at an initial position and velocity of the satellite that are guessed,
and that gradual improvements are made to the computed trajectory so that it comes closer
to the true trajectory. The consequence is that the partial derivatives are obtained by solving
eq. (10.37) also depend on the initial guess orbit, and that those partial derivatives require an
iterative improvement to bring them closer to the true trajectory of the vehicle.

141
Bayesian parameter estimation
In addition, the parameter estimation problem that originates for satellite problems is poorly
conditioned so that constraint equations as discussed in section 8.8.4 are included in the pa-
rameter estimation procedure. For this reason the Bayesian parameter estimation method as in
eqn. (8.64) is normally used to be able to find a solution for the estimation problem. One reason
may be that the control parameters are not well determined in a least squares approach because
a limited time span was used in the variational problem, another reason may be that initial
state vector control parameters are correlated with dynamical parameters. Bayesian parameter
estimation methods turn out to be useful especially when the involved constraint equations are
f
based on external information that is available for in eq. (10.37). Also, constraint equa-
tions may need to be formulated for the initial state vector, for instance because the argument
of perigee and the true anomaly become linearly correlated for near circular trajectories of the
satellite.

Partitioning of parameters
Normally variational problems for satellite orbit determination are confined to a limited integra-
tion time, the reason is usually tractability because of the used numerical integration techniques.
A drawback is that certain dynamic parameters are estimated again while they are specific for a
defined integration interval. In this case we speak about arc parameters where an arc is defined
as the trajectory that is confined in duration, for satellites at altitudes below 1000km the arc
length is usually between one day up to a week. If observations to a satellite are combined from
several arcs then two type of parameters may occur, namely those parameters in the problem
that are specific for an arc, and all other parameters that are spanning multiple arcs. Blocking
methods as introduced in section 8.9.3 are often used 1) to solve the arc parameters by making
use of eq. (8.72), and to 2) continue with the reduced normal equations that are obtained with
eq. (8.73). This procedure is known as backsubstitution.
The discussion also shows that users may significantly benefit from well calibrated force mod-
els to avoid that partitioning of parameters is necessary in their orbit determination procedure.
Nowadays this is the standard practice in precision orbit determination, in most cases use is
made of the procedures recommended for instance by the IERS, cf. [44] specifically to avoid that
a precision orbit determination task automatically becomes a task where many parameters in
planetary physics models need to be determined over and over again.
Batch least squares parameter estimation has significant advantages but inefficient partition-
ing of the involved parameters may easily lead to a significant burden to solve the parameter
estimation task. One example is the modeling of drag and solar radiation pressure parameters in
an orbit determination problem. The batch least squares problem would at best cause a patched
polynomial approach whereby drag model parameters Cd or solar radiation pressure constant
Cr are estimated. The arc length for a satellite at 500km may in this case be a week, but, the
variability of the forcing as a result of the thermosphere and the solar radiation is too large to
define one version of the Cd parameter and one version of the Cr parameter for the entire week.
In reality Cd parameters are then estimated in a 3 hourly patches whereby each patch assumes a
constant value for Cd . For Cr the forcing is smaller, longer patches are then used, typically they
are 12 hours or more in length. As a result several hundred parameters are introduced in the
variational problem, which can be inefficient. Yet the essence is that there are no observation

142
equations for range or range rate with overlapping polynomial patches. The result is that the
design matrix will become sparse, in fact, it will become diagonal dominant because no inner
products between columns of the design matrix will be formed. The alternatives are then to
continue with sparse matrix techniques as introduced in section 8.9.2, but another method is
to use the outcome of the variational method in a more efficient way by means of sequential
parameter estimation.

10.3.5 Sequential parameter estimation

One possibility to increase the efficiency of a batch least squares method is to directly make use
of the fact that the problem is partitioned, and that information of a previous patch is used to
assist the parameter estimation procedure in the next patch. For this purpose we will define a
state vector xj which is defined at epoch tj . In addition there is a covariance matrix Pj that
is associated with xj . We will advance from epoch tj to the next epoch tk and we have at
our disposal a state transition matrix (tk , tj ) or short kj . For satellite orbit determination
problems the state vector is obtained from a variational problem as in eqns. (10.36), but the
technique of sequential parameter estimation is general, and it may be applied to other dynamical
systems as well.
The Kalman filter1 as we know it today is a well known method for sequential parameter
estimation. The method considers a state vector xj that comes with a covariance matrix Pj
which as both defined at epoch tj . The Kalman filter algorithm consists of two distinct steps:

1. There is a propagation step whereby the state vector xj and it associated covariance matrix
Pj are transformed into the combination xk and Pk . This means that information at epoch
tj is propagated to epoch tk which is accomplished by straightforward application of the
relation that xk = kj xj . Since the relation is linear we can also transform the covariance
matrix, simply by Pk = kj Pj tkj . This is a consequence of the properties discussed in
section 8.4.3.

2. Once we arrive at epoch tk with the propagated information from tj the following situations
may occur. Either there are no observations, in which case our only option is to continue
propagating to future epochs. But the other possibility is that there is observational data
in the form of a observation vector y k with a corresponding covariance matrix Rk . This
leads to constrained parameter estimation as discussed in section 8.8.4, it is implemented
in the so-called update of the Kalman filter.

Kalman filter implementation

As was explained before, at epoch tj we have the state vector xj and its covariance matrix
Pj . For the first epoch in the Kalman filter any value can be selected, for instance, that the
initial state vector is zero and that the covariance matrix is diagonal with exaggerated variances
on the main diagonal. This will not matter because the Kalman filter will rapidly update the
state-vector and the covariance as new observational data comes along. The propagation from
epoch tj to tk yields:
xk = kj xj
(10.38)
Pk = kj Pj tkj
1
Named after Rudolf Emil Kalman who was born in 1930 in Budapest

143
at epoch k we assume that observation data is available in the form of data vector y k , also,
there is an information matrix Hk that relates y k to parameters in the state vector xk . The
observation equations at epoch tk are:

y k = Hk xk +  with Rk = E( t ) (10.39)

The update step at epoch tk will now combine two sources of information, namely, the propagated
information from epoch tj and the observation information at epoch tk . The update step at epoch
k is the Bayesian parameter estimation algorithm discussed in section 8.8.4, as a result:

1
x = Hkt Rk1 Hk + Pk1 Hkt Rk1 y k



1 (10.40)
Pk = Hkt Rk1 Hk + Pk1

where the hat symbols above the vector and matrix on the left hand side indicate that this
information follows from a least squares optimization. These equations are close to what we find
in literature as the Kalman filter.

Remark 1: add the predicted state vector xk in the update step

In equation (10.39) we did not make use of the fact that there is prior information for xk which
we could have used, also, in step (10.40) we ignored prior information in the form of Pk . Kalman
therefore re-formulated the problem. Let us first begin with the assumption that xk is known
at epoch k and that it is used to de-bias the observation equations. Kalman used a gain matrix
Kk that appears in the update equations where the observations at epoch tk are corrected for
the predicted observations that follow from Hk xk .

xk = xk + Kk (y k Hk xk ) (10.41)

where the gain matrix follows from the properties of the Penrose-Moore pseudo inverse as dis-
cussed in section 8.7.2: 1
Kk = Pk Hkt Hk Pk Hkt + Rk

(10.42)
Covariance propagation is applied to eq. (10.41), and this results in the update equation for the
covariance matrix at epoch tk :
Pk = [I Kk Hk ] Pk (10.43)

Remark 2: assume that the dynamics in kj is not complete

In eq (10.38) we assumed that the propagation of the filter is perfect, and that all dynamical
effects are known and represented in kj . In reality this is not the case so that there is a
need to define system noise. Propagation of the state vector and the covariance matrix often
lead, on the long run, to situations where the propagated covariance matrix weighs too heavily
on the Kalman filter update equation. As a result the Kalman filter becomes insensitive for
new observation information because it thinks that the state vector and the covariance matrix
propagated from the previous step are too accurate. To compensate for this situation a so-called
state noise compensation (SNC) algorithm is devised, and it is implemented by assuming that
the propagation step takes another shape:

xk = kj xj + kj uk (10.44)

144
where uk contains system noise for which we assume that:

E(u(t)) = 0 (10.45)

and
E(u(t)ut ( )) = Q(t)(t ) (10.46)
The consequence for the propagated covariance is that:

Pk = kj Pj tkj + kj Qk tkj (10.47)

Details about this algorithm are discussed in section 4.9 in [63]. The consequence of considering
process noise in the SNC algorithm is that the optimism in Pk is reduced, and that the Kalman
filter does not stall meaning that it becomes inert to new observation information added at tk .

10.3.6 Toy Kalman filter without process noise

This problem assumes that the state vector consists of two variables, namely T (temperature)
and its derivative to time dT /dt and that there are updates once every 60 seconds in the form
of temperature observations. For the toy problem we want to demonstrate that the Kalman
filter will reach a steady state and that it becomes inert since process noise is not part of the
algorithm. We start with the assumption that:
 
1 t
kj = (10.48)
0 1

where t = 60 and that:  

T
xk = (10.49)
dT /dt
with T representing temperature. Furthermore we assume:
 
1000 0
P0 = (10.50)
0 1000

and  
0
x0 = (10.51)
0
For the design matrix we assume that there is an observation batch where T is observed at epoch
tk and that all elements in the observation batch are uncorrelated.

1 0
Hk = ... ... (10.52)

1 0

1
Rk =
..
(10.53)
.
1

145
 Figure 10.4: Example of Kalman filter without process noise, = 0.

The observation equations at epoch tk consist of n independent samples, and this is formulated
as follows:
y1 1 0  
.. .. .. T
. = . . + tk (10.54)
dT /dt t
k
yn t 1 0
k

This method was implemented in matlab where we fed a synthetic dataset to the filter, and
it produced the results shown in figure 10.4. The top left part in figure 10.4 shows the filter
output in red, and the simulated observations in blue. The other panes show, bottom left: the
dT/dt variable in red predicted by the filter. In the right side panes one finds the corresponding
standard deviations (square roots of the diagonal elements of the predicted covariance matrix)
that follow from the filter. It can be seen that the predicted temperatures tend to follow the
observations, but at some point in time (around epoch 4000, the units are in multiples of 60
seconds) new information does not change the filter output. At this point we say that the
Kalman filter has approached a steady state but we also see that it has become inert. The latter
is a consequence of the fact that we ignored to formulate process noise which typically avoids
such situations.
In the above example we see that the Kalman filter output smooths the observations, in
the beginning the filter is close to the T observations, but as we go along it becomes more and
more resistent to any new input. This is what we see in the upper left pane of figure 10.4, in
the upper left figure we see that the standard deviation peters out (asymptotically reaches) at
approximately 0.1, for the derivative in the lower right this is even worse, so the Kalman filter
has in this case really become insensitive to new observations.
The other benefit that we see from using the Kalman filter is that derivatives of T are freely
produced, there is no need to create for instance polynomials that span a number of epochs,
and also, there is no need to buffer observations over an extended period of time. Editing of

146
 Figure 10.5: Example of Kalman filter with process noise, = 108 .

bad observations may be performed on the fly, that is, as new observations come along then we
have the possibility to calculate tk and we can verify whether these residuals correspond to for
instance Rk . In this sense there are various possibilities to adaptively modify elements of the
Kalman filter on the fly, or to reject bad observations. And these possibilities would not exist if
we used a batch least squares method.

10.3.7 Toy Kalman filter with process noise

We leverage on the experience gained version 1 of the Kalman filter and we modify the algorithm
in the following way:
xk = kj xj + kj uk (10.55)
where kj = I and E(uk utk ) = Qk = I. When the filter dynamics is changed we also affect the
update step of the Kalman filter, and this goes along the lines discussed in eq. (10.41) to (10.43).
We can now run simulations where values are changed to inspect how the predictions made by
the Kalman filter are influenced. Figure 10.5 is an example where weve set = 108 , effectively
this adds a little bit of process noise to the algorithm so that we assume less weight for the
predictions in the update step. By adding process noise we accomplish that the filter becomes
less sensitive to the predicted dynamics, and that it becomes more responsive to observations
(read data) added at each update step.

10.3.8 Characterization of process noise

In reality adding process noise is not as simple as is shown here, because oftentimes use is
made of specific models to characterize the behavior of the noise spectrum. The behavior of
instruments and in particular, specific components within those instruments such as oscillators,
gyroscopes, star cameras and phase discriminators help to specify process noise in a Kalman

147
filter. Reference [53] gives a summary of possible spectral density models that can be used for
quartz-crystal oscillators. It mentions that spectral density can be modeled as Sy (f ) = h f
where f is the frequency with h characterizing the noise level. Slope parameter applies the
type of noise where can vary between -2 and +2. Table I in [53] summarizes spectral laws for
to characterize different regimes.

= 2 : Random walk frequency noise

= 1 : Flicker frequency noise

= 0 : White frequency noise

= 1 : Flicker phase noise

= 2 : White phase noise

In [53] it is recommended that spectral analysis techniques should be consulted to determine

h f within different regimes. A technique that is often used for this purpose considers the
definition of Allan variances.

10.3.9 Allan variance analysis

In chapter 8 we presented a general approach for estimating parameters where we encountered
expectancy operators including definitions of averages (means), medians, variances and estima-
tion procedures. The definition of Allan variances is an extension of what has been presented
in that chapter, and it is used to classify the spectral density of a noise variance model.
Allan variances follow from a measurement series that is regularly spaced, the data consists
for instance of successive frequency readings of a clock oscillator. All clocks have a circuit
called an oscillator that generates a high number of oscillations per second, those oscillations
are counted and the outcome is translated into a counter reading that we attach to define an
epoch. Divide the counter reading over the reference number on the crystal used in the oscillator,
and you have a measure for the second. Another possibility of the read-out the clock counter of
a guest clock that is observed with a more accurate observer clock. It is this second example
that we will use to gather a frequency dataset for the Allan variance calculation.
In our clock experiment we could gather frequency values every 10 seconds (this is the so-
called sampling interval) and we could continue this procedure between 0 and 10000s. Experience
tells that the obtained variance (and mean) of the observed frequencies in the dataset will depend
on 1) the length of the data record, and 2) the sample interval. The choice of dataset length
and sampling interval are arbitrary, the consequence of this is that we do not get a complete
overview of what we can be expected from the variance of a clock oscillator.
Allan variances partially solve this problem, because they do take into account a measure
for the variance as a function of the sample interval from a series of frequency measurements
collected by the observer clock from the guest clock oscillator. There are two sort of datasets that
we could process, namely datasets that contain the phase of the guest clock and datasets that
contain the frequency of the guest clock. In the phase dataset we collected counter readings at
the sampling interval, in the frequency dataset we divide the phase values of the guest clock
over the sampling interval defined by the observer clock.

148
 Allan variances are derived from so-called M-sample variances, if the dataset contains phase
measurements x(t) then:
#2
1 
"M 1
1 MX x(iT + ) x(iT )
2
1 X x(iT + ) x(iT )
y2 (M, T, ) = (10.56)
M 1 M
i=0 i=0

where M denotes the number of samples in a data record as described before, T is the sampling
interval time, is a parameter that we have introduced ourselves, it is the integration time of
the frequency estimate. The analogy with a dataset of frequency measurements is:
(M 1 "M 1 #)
2 1 X
2 1 X 2
y (M, T, ) = yi yi (10.57)
M 1 M
i=0 i=0

Allan variances are now defined as y2 ( ) = y2 (2, , ) where hi is short for an averaging



integral. Allan variances are a measure for the sensitivity of clocks frequency variance as a
function of the chosen integration time .
An example of various Allan variances is taken from [27] where different clock design are
presented, cf. figure 10.6. In this figure the Allan deviations (square root of the variance) of
the relative frequency error (f /f ) of various clocks are shown with the logarithmic values of
along the horizontal axis. The discussion clearly shows that the variance model of a clock
oscillator has three different regimes, namely flicker frequency noise and random frequency noise
from 0s up to a specified integration time (the reason is that the phase of the oscillator is
sampled, and that the sampling error affects the measurement), then there is a white frequency
noise floor (this is apparently the best you can get out of an oscillator), and finally there is white
or random walk phase noise when the integration time extends (in this case we get to see the
long term scintillations in the frequency that build up in time, it may exhibit a variety of effects
affecting the performance of the oscillator). Notice also how the shape of each curve depends
on the hardware used in the clocks oscillator, Cs stands for Cesium, Rb stands for Rubidium,
X-tal stands for quartz crystal, and H stands for a Hydrogen maser.
The consequence of the theoretical clock model is that the definition of noise comes from the
regimes of the Allan variances described in [53]. Allan variances are not only defined for clock
oscillators which form the basis of many instruments, the same procedure may be implemented
for all components within a measurement system. One of the possible applications of Allan
variances is to specify Kalman filter parameters as is discussed in [64].

149
Figure 10.6: Allan clock variances of clocks, figure comes from [27]

150
Chapter 11

Three body problem

The two-body problem in celestial mechanics considers a satellite orbiting a planet (or the Sun)
whose mass is far greater than that of the satellite. The motions are usually circular or elliptical;
but can also be parabolic or hyperbolic. For circular and elliptical trajectories the orbital period
depends on the mass of the planet and the semi-major axis of the satellite. In the three-body
problem we add another planet to this configuration whose mass is about the same order as
that of the first planet. There are no easy solutions for the three body problem, but there
are approximated solutions for the restricted three-body problem. The main result is shown in
figure 11.4 where we have two planets P and Q and 5 Lagrange points, wikipedia tells us that
Euler discovered L1 to L3 a few years before Lagrange found L4 and L5 in 1772.

11.1 The restricted three-body problem

In figure 11.4 we assume that planet P with mass mp is located at (0, dp ) and that planet Q
with mass mq is at (0, +dq ). The system rotates with a constant angular speed n about a center
of mass of the system (also known as the barycenter) which is the meeting point of all dashed
lines in figure 11.4. The angular rate of both planets will depend on the sum of their masses
and the sum of dp and dq . These distances follow in turn from the masses mp and mq .
For the restricted three-body problem we demand that the sum of the centrifugal and the
gravitational contributions of the acceleration balance one another for each planet. Therefore:
1 Gmp mq q
n2 dp = 2
= (11.1)
mp (dp + dq ) (dp + dq )2
1 Gmp mq p
n2 dq = = (11.2)
mq (dp + dq )2 (dp + dq )2

which leads to the mean motion n for this problem:

p + q
n2 = (11.3)
(dp + dq )3

Compared to the two-body problem we can simply replace the gravitation constant of the Sun
(or main planet) by the sum of both gravitational constants of both planets in the three-body
problem. An essential step in the restricted problem is to assume that the mean motion n is
constant. We derive the equations of motion in this system by introducing a transformation for

151
Figure 11.1: (a) The left figure shows the orbit for a small (red) particle and much heavier planet
in the blue orbit with the Sun in the origin. The Sun is 1000 times heavier than the planet,
and the planet is far more heavy than the particle. (b) The right figure shows the orbit of the
particle represented in the rotated system, the planet is always located at (1000,0), and we used
for in eq.(11.4) a suitable value derived from the argument of latitude of the planet.

in a rotating system. The purpose of the transformation R3 is to transform local coordinates

from the rotating system to inertial positions in the inertial system x:

x = R3 () (11.4)

where (t) = n.(t t0 ) is a linear rotation angle and R3 the rotation matrix. In the following
section we will introduce a rotated coordinate system with an example.

11.2 Two bodies orbiting the Sun

Figure 11.1 shows an example of two circular orbits which we will use to illustrate the concept
of a rotated coordinate system. In figure 11.1a the orbit of the planet starts at (1000,0), and the
motion is counterclockwise. At the same time the particle in the red orbit starts at (-2000,0)
and it moves in the same direction but its speed is clearly slower. (With Keplers laws you
can verify the velocity ratio of both trajectories.) The top-view is often chosen to illustrate
an orbit problem, but, for the Lagrange problem it is more convenient to consider a rotating
system, where we simply apply equation (11.4) to all coordinates shown in figure 11.1a. The
transformation is implemented in such a way that the minor planet will be fixed at point (1000,0)
and the Sun in (-1,0) after transformation. The rotation is now applied to the red orbit, and
under this assumption we get figure 11.1b which displays the motion of the particle relative to
the planet and the Sun.
As you can see, in figure 11.1 there is hardly any interaction between the particle and the
planet. In fact, all motions seem to be determined by the Sun in this configuration. However,
this will soon change once we lower the orbit of the particle. Figure 11.2a shows like figure 11.1
a top view of the solar system for a particle in an orbit whose radius is just 3% larger than that
of the planet. Over time we will see that the particle slightly lags the planet, but the difference

152
Figure 11.2: (a) The left figure is a top view of the (red) particle orbit, and (b) shows the rotated
motion relative to the planet. The radius of the particle is now 3% larger than that of the planet
which is at (1000,0)

in orbital speed is small. This can be seen in figure 11.2b where the particle orbit is shown in a
rotated view relative to the planet at (1000,0). As you can see, in figure 11.2b the particle stays
close to (-1030,0) and it only gradually moves northward after one orbit of the planet about the
Sun. At some time one might expect the particle to approach the planet at (1000,0) so that
the gravitational influence of the planet on the particle will become stronger. Lets see what
happens, in figure 11.3 we extended the time window in the numerical orbit integration process;
we are surprised to see that the particle didnt reach the planet at (1000,0), instead the particle
turns around and it loops back close to its starting position to approach the planet from the
other side along L5 , for this reason it is called a horseshoe orbit. The existence of Lagrange
points starts here where we ask the following questions: 1) Why do we see motions like shown
in figure 11.3 and 2) are other particle motions possible? In order to understand this problem
we rewrite the equations of motion in a rotated system. A little later we will also look at the
energy in the system.

11.3 Accelerations in a rotating system

In order to explain what we observed in section 11.2 we need to apply a rotation to the equations
of motion so that apparent forces (Coriolis effects and centrifugal forces) will appear in the
problem. Lets first start with the inertial system, which is equivalent to the non-rotated top
view of the solar system. In this system the equations of motion are:

x1
p q
x = x2 = (x xp ) + (x xq ) (11.5)
|x xp |3 |x xq |3
x3

153
Figure 11.3: Particle orbit extended over time and represented in the rotated system

154
where the vectors xp and xq model the Sun and the planet respectively while x is the position
of the particle. It is relatively easy to show that:

1 2n2 n2 1

p q
x = 2 + 2n1 n2 2 = 3
( p ) + ( q ) (11.6)
| p | | q |3
3

where = (1 , 2 , 3 ) is the position of the satellite and where p and q are the positions of P
and Q in the rotated frame, also called the frame. Eq. (11.6) was used to generate the plots
in the rotated frame. We used a Matlab procedure to solve a system of first-order differential
equations (ODEs). Both ode45 and ode115 work fine as long as you set the relative and
absolute tolerances on the numerical integrator error to approximately 1012 . To accomplish
the latter you use matlabs odeset routine. Furthermore you should rewrite the second order
equations of motion shown here as a system of first-order ODEs where you provide ode45 and
ode115 a link to your function that calculates y = F (t, y) where y is a state vector and t is
time. Matlab has a great help function and documentation that clarify all its features, and it
is freely available to all TU Delft students.
To understand the rotated frame results that we found for the particle we will now plot the
length of . If we exclude the velocity of the particle (which would introduce a Coriolis effect)
and if we constrain the motion to a plane (x3 = 3 = 0) then we obtain experienced by a
particle in a rotating frame. This assumption results in:
1 + dp 1 dq
1 = p 3
q + n2 1 (11.7)
| p | | q |3
2 2
2 = p q + n2 2 (11.8)
| p |3 | q |3

The length of the acceleration vector || can now be plotted as a function of the position in the
frame. This is done in figure 11.5 where we have assumed a hypothetical configuration with
dp = 1, dq = 10, p = 10 and q = 1.
The position of Lagrange points are shown in 11.4, these points will appear within the
gravity-wells which are located in the blue regions in figure 11.5. The first well is the C
shaped horse-shoe where L3 , L4 and L5 can be found in the white exclusion zones. The
second well is between P and Q, the third well is located behind Q when facing it from P .
In figure 11.5 we have ignored large accelerations in the neighborhood of P and Q where the
local gravitational effect is dominating. Furthermore we ignored to plot || in the outer region.
In figure 11.4 we indicate the corresponding Lagrangian points L1 to L5 where a satellite would
not experience any residual acceleration because || = 0. In these regions there is a balance
between gravitational and centrifugal accelerations so that the netto acceleration is zero.

11.4 Jacobi constant

In [11] you will not exactly find a plot like shown in figure 11.5, instead you will find a Jacobi
constant Cj which is defined as:

p q 1 1
Cj = + n2 (x21 + x22 ) + (x1 2 + x2 2 + x3 2 ) (11.9)
|x xp | |x xq | 2 2

155
 L4

P Q

L3 L1 L2

L5

Figure 11.4: The restricted three body problem; positions of Lagrangian points are indicated by
open circles, P and Q are planets

The Jacobi constant gives us the total energy of a particle in the three-body problem. The
discussion relates to the existence of so-called Hill surfaces within which a particle can remain
as long as the total energy does not exceed a limit. The latter would be possible if we gave a
particle too much velocity for instance. Figure 11.6 shows the Jacobi constant for a hypothetic
case, we increase the mass ratio so that the gravity wells stand out.
Figure 11.7 shows a perturbed particle orbit that started near L3 , and figure 11.8 is an
example how a perturbed motion is constrained near L4 . In [11] you will find examples in our
Solar system that look like horseshoe or tadpole orbits. L4 and L5 are regions where one can
find Trojan asteroids in the Sun Jupiter system, and the motion of the moons Epimetheus and
Janus in the Saturnian system closely resembles the horseshoe motion shown in figure 11.7, see
also figure 11.9.

11.5 Position Lagrange points

Lagrange points L1 L2 and L3 should appear on the line connecting the planets, and L4 and
L5 appear at angles of 60 relative to P to this line (it is really P and not the origin, please
check this yourself). In this configuration all accelerations will cancel in L4 and L5 . A little
more effort is required to locate the other Lagrangian points. From eq. (11.6) we conclude that
a particle can only move on the line connecting L1 L2 and L3 when 2 = 0 so that 2 = 0. To
locate L1 L2 and L3 we must solve s in:
p q
n2 s 3
(s + dp ) (s dq ) = 0 (11.10)
|s + dp | |s dq |3

156
 20

15

10

10

15

20
25 20 15 10 5 0 5 10 15 20 25

Figure 11.5: Colors indicate the length of the local acceleration vector in the rotating coordinate
system (blue colors indicate shorter values than red ones). This configuration is computed for
p = 10, q = 1 and correspondingly p = (1, 0) and q = (0, 10).

157
Figure 11.6: Jacobi constant plot for a synthetic case where p = 20 and q = 1 which is far
greater than any mass ratio that we have in the solar system, e.g. the Earth Moon mass ratio
is 80.

158
Figure 11.7: Example of a horseshoe orbit that started in L3 , the orbit is confined within
the gravity well around L5 L3 and L4 . Over time this type of orbit should be considered as
marginally stable.

159
Figure 11.8: Example of a tadpole orbit that started in L4 , in this case the orbit is confined to
a narrow region around L4 (or L5 ). This type of orbit should be considered as stable.

Figure 11.9: Epimetheus and Janus orbiting Saturn, credits Cassini project JPL/ESA.

160
 p /q s0 at L1 s0 at L2 s0 at L3
10 0.795450 1.256083 -1.037836
100 0.848624 1.146320 -1.004125
1000 0.931310 1.069893 -1.000416
10000 0.968066 1.032424 -1.000042

Table 11.1: Positions of Lagrange points L1 to L3 , here s0 = s

dp +dq

where s is an ordinate along the connection line counted from the center of mass of the system. To
determine the roots of equation (11.10) I recommend the use of maple or a numerical root finding
routine (hint: familiarize yourself with either maple or matlab). The roots of equation (11.10)
depend on the ratio of p and q ; some examples are shown in table 11.1 where it is assumed
that mq = 1 and G = 1. Lagrangian points L4 and L5 are easy to find. In this case we assume
that || 1 and | p | 1 and | q | 1 which is valid when p >> q so that = 0 which
directly follows from eq. (11.6). This situation is only possible for a configuration where the
Lagrange points are located on the top of a triangle with sides of length 1, see also figures 11.4
and 11.5.

11.6 Stability conditions in Lagrange points

.
One may wonder whether motions that start near the Lagrangian points will or will not
remain stable for an extended period of time. This discussion is further worked out in the
astrodynamics I lecture notes where the local eigenvalue problem is considered. As a result
one can investigate whether the solutions will either oscillate and whether these oscillations will
reduce over time.The approach is to linearize equation (11.6) at the Lagrangian points and to
investigate the eigenvalues of the coefficient matrix. The discussion results in obvious instability
at Lagrangian points L1 and L2 because of the saddle point in the Jacobi constant function. In
L3 one can have marginally stable motions but it may be expected that the particle orbit can
eventually reach an unstable regime, finally, for L4 and L5 the orbits are stable over time.

11.7 Exercise
The restricted three body problem considers with two planets, P1 and P2 and a small particle
P3 where all bodies should be considered as point masses. This set-up resulted in the definition
of Lagrangian points L1 to L5 . In the following questions we ask you to look into the dynamics
of a this problem:

1. Formulate the equations of motion of all bodies in this system in an inertial coordinate
system.

2. Assume that P2 orbits P1 in a circular trajectory. Formulate the equations of motion for
P3 in a rotating system where the x-axis is aligned with P1 and P2 and where the z-axis
is aligned with the angular momentum vector of the system.

161
3. Write a matlab script where you implement the equations of motion under b) for P3, and
modify the initial conditions until youve found an orbit that stays in the gravity well
around Lagrangian point L3.

162
Chapter 12

Co-rotating orbit dynamics

G.W. Hill (1838-1914) considered equations of motions in a coordinate system that rotates at
a uniform rate. Although the method seems theoretical it turns out that this approach is very
useful for approximating the effect of perturbing forces on satellites but also planets in the solar
system. Also, the Hill problem makes it relatively easy to identify the approximate positions
of Lagrangian points L1 and L2 in the three body problem. To derive the Hill equations we
consider an inertial system x where the x and y axis appear within the orbital plane and where
the z-axis is pointing in the direction of the angular momentum vector. Consider also the
rotating system with the -axis coinciding with the z axis and where the -axis is pointing at
the satellite. The situation sketch is shown in figure 12.1. The relation between both systems is

Figure 12.1: The x and the system.

163
as follows:
x cos sin 0
x = R() y = sin cos 0 (12.1)
z 0 0 1
and
(t) = 0 + t = 0 + nt (12.2)
where n is constant. The second-order derivative of x with respect to time is:
x = R + 2R + R (12.3)
so that:
2n n2

x
y = R() + 2n n2 (12.4)
z
The same is true for the gradient of the potential:

V /x V /
V /y = R() V / (12.5)
V /z V /
It is relatively easy to show that R() contains the partial derivatives and that (12.5) follows from
the chain rule. The equations of motion in the rotating system are obtained from eq. (12.4)
and eq.(12.5):
V
2n n2 =

V
+ 2n n2 = (12.6)

V
=

The next step is to express the potential in the frame at the position of the satellite. This is
accomplished by linearizing the expression V = U + T with U = /r at the nominal orbit, cf.
= r, = 0 and = 0 where T is referred to as the disturbing potential. Furthermore u, v
en w denote small displacements in the frame. (u = , v = , w = .) The linearized
gradient in the frame at the true position of the satellite (cf. = r + u, = v en = w) is:
V T
= 2
+ 2 3u + +
r r u
V T
= 3v + + (12.7)
r v
V T
= 3w + +
r w
From equations (12.6) to (12.7) we find:
T
u 2nv 3n2 u =
u
T
v + 2nu = (12.8)
v
T
w + n2 w =
w

164
and these equations are known as the Hill equations.

12.1 Solution of the Hill equations

An interesting property of the Hill equations is that analytical solutions exist provided that n
is constant. In this case eq. (12.8) becomes:
u = F u + g (12.9)
where F does not depend on time and where g does depend on time. The homogeneous and the
particular solutions are discussed in the following sections.

12.1.1 Homogeneous part

In this case g = 0 and to obtain a solution we decompose F in eigenvalues while the eigen-
vectors appear in the columns of the Q matrix:
F = QQt (12.10)
The homogenous solution becomes:
u(t) = Qe.(tt0 ) Qt u(t0 ) (12.11)
To demonstrate that this is a solution we consider a Taylor expansion of u(t) :
1 1
u(t) = u(t0 ) + u(t0 )(t t0 ) + u(t0 )(t t0 )2 + . . . + u(n) (t0 )(t t0 )n (12.12)
2 n!
with:
u = F u
u = F u + F u = F.F.u = F 2 u
..
.
nu
= F nu
dtn
so that:
1 1
u(t) = (I + F.(t t0 ) + F 2 .(t t0 )2 + . . . + F n (t t0 )n + . . .)u(t0 )
2 n!
from which it follows that:
u(t) = eF.(tt0 ) u(t0 ).
If F = QQt then:
t .(tt )
u(t) = eQQ 0
u(t0 ),
and as a result:
t .(tt )
eQQ 0
= Qe.(tt0 ) Qt (12.13)
This shows that eq.(12.11) is a homogeneous solution since Q is orthonormal. Since .(t t0 )
is diagonal the term e.(tt0 ) is easily obtained, for real i on the diagonal of we obtain
exponential expressions that either decay for negative eigenvalues or that grow for positive
eigenvalues and for complex eigenvalues i we find oscillating solutions.

165
12.1.2 Particular solution
For the homogeneous part we found:
y0 + a y = 0 (12.14)
where
(x) = ceax
is a solution. It can be shown by a substitution of:

0 (x) = c aeax (12.15)

in eq. (12.14). For the particular solution we need:

y 0 + ay = b(x) (12.16)

for which it can be shown that:

(x) = ceax + eax B(x) (12.17)

with Z x
B(x) = eat b(t) dt (12.18)
x0

This can be shown by substitution of 0 (x) and B 0 (x) in eq. (12.16). We remind that Eq. (12.18)
is known as the Laplace transform of b(x) and that Laplace transforms of most functions are
known. To demonstrate that the same technique can be used for a system of differential equations
we consider the system:
y + Ay = b(x) (12.19)
where we decompose A as:
A = QQT
and where we pre-multiply with Qt :

QT y + QT y = QT b(x)

If z = QT y then:
z + z = c(x)
so that we obtain a decoupled system. In this case the solution is:

(x) = dex + ex C(x) (12.20)

Z x
C(x) = et c(t) dt (12.21)
x0

where
y(t) = Q (x) (12.22)
is a solution of (12.19).

12.2 Characteristic solutions

Characteristic solutions are known for all possible variations of the Hill equations.

166
12.2.1 Homogeneous solution
The homogeneous system takes the following form:
u 2nv 3n2 u = 0
v + 2nu = 0 (12.23)
2
w + n w = 0
and the characteristic solution is:
u(t) = au cos nt + bu sin nt + cu
v(t) = av cos nt + bv sin nt + cv + dv t (12.24)
w(t) = aw cos nt + bw sin nt
where the integration constants on the right hand side depend on the initial conditions of the
problem.

12.2.2 Particular solution

The particular system takes the following form:
u 2nv 3n2 u = Pu cos t + Qu sin t
v + 2nu = Pv cos t + Qv sin t (12.25)
2
w + n w = Pw cos t + Qw sin t
and the characteristic solution is:
u(t) = au cos t + bu sin t
v(t) = av cos t + bv sin t (12.26)
w(t) = aw cos t + bw sin t
where the integration constants on the right side depend on the the forcing parameters in
eq. (12.25).

12.2.3 Particular resonant solution

The particular resonant system is an exception that we need when n = or when = 0 in
the standard particular solution:
u 2nv 3n2 u = Pu cos nt + Qu sin nt + Ru
v + 2nu = Pv cos nt + Qv sin nt + Rv (12.27)
2
w + n w = Pw cos nt + Qw sin nt + Rw
and the characteristic solution is:
u(t) = (a0u + a1u t) cos nt + (b0u + b1u t) sin nt + c0u + c1u t
v(t) = (a0v + a1v t) cos nt + (b0v + b1v t) sin nt + c0v + c1v t + c2v t2 (12.28)
w(t) = (a0w + a1w t) cos nt + (b0w + b1w t) sin nt + c0w
where the integration constants on the right side depend on the forcing parameters in eq. (12.27).

167
12.3 Exercises
1. Solve the integration constants of the homogeneous solution of the Hill equations, an
engineering application for this homogeneous solution may be the rendezvous problem.

2. Show that the particular non-resonant u(t) equation becomes:

2nQv + Pu 2nPv + Qu
u(t) = 2 2
cos t + sin t (12.29)
(n ) (n2 2 )

It shows that the response of the system is linear, that is, if you apply a perturbing
acceleration on a particle at frequency then orbit perturbations appear at that same
frequency . The latter equation can be used to demonstrate long term stability in the
Solar system.

168
Chapter 13

Hill sphere and Roche limit

13.1 Hill sphere

The Hill sphere of a planet is defined as the approximate limit of the gravitational influence
of the planet in orbit about the Sun. The definition may be extended to a moonlet inside a
ring of a planet in which case the spheres radius follows from the presence of the moonlet near
the planet. Within the Hill sphere of a planet we can find moons, outside the Hill sphere a
moon can not stay near the planet because the gravitational effect of the Sun dominates. At the
Hill sphere both accelerations are in balance, so that the radius should extend to Lagrangian
points L1 and L2. The Hill sphere is therefore the limit between a two and three body problem
mechanics. According to [11] the approximate extent of the Hill sphere is:
 1/3
mq
Rh a (13.1)
3(mp + mq )
and the question is now, why is this the case? To demonstrate this relation we consider the Hill
equations as shown in eq. (12.8). Due to the geometry we only need to consider the u equation,
the second and third component are not relevant, so that v = 0 and w = 0. For this problem
the u equation becomes:
G.mq G(mp + mq )
u = 2 = 3n2 u = 3 Rh (13.2)
Rh a3
where we used the definition of n from the three-body problem. As a result we get:
 1/3
3 mq 3 mq
Rh = a Rh = a (13.3)
3(mp + mq ) 3(mp + mq )
At the same time, this relation may be used to approximate the location of L1 and L2, a
numerical algorithm such as Newton Raphson procedure may then continue to optimize the
roots of L1 and L2 with eq. (11.10). For L3 the Newton Raphson algorithm can start at s = 1.
Note however that eq. (13.3) should be rescaled confirm the definition of s in eq. (11.10).

13.2 Roche limit

In [11] background information is provided on the definition of the Roche limit, the essence of
the problem is to find the minimal distance between a planet and a satellite so that the tidal

169
acceleration at and the binding acceleration ab balance at the satellite. Thus:

at = ab (13.4)

In the following we will assume that:

The planet has a mass Mp and a radius rp , its gravitational constant is p and the its
density is p

The satellite has a mass Ms and a radius rs , its gravitational constant is s and the its
density is s

The separation distance between planet and satellite is called d

So far we have not said where the balance holds and how the binding or tidal acceleration should
be calculated. In fact, this depends on how you exactly define the problem. The straightforward
method is to assume that the satellite is at distance d and that the balance holds at its surface.
In this case you get, see also [11] and chapter 14 for more detail:
3p s
3
rs = 2 (13.5)
d rs
where the left hand side is obtained via a Tayler series approximation of the gravitational
attraction at the satellites center times the linearization distance rs . The right hand side is the
opposite acceleration at the satellites surface. We arrive at the expression:
p 3
d3 = 3 r (13.6)
s s
where the ratio of the gravitational constants of planet and satellite can be reduced to:

p p rp3
= (13.7)
s s rs3
so that the Roche limit becomes:
 1/3
p
d = 1.44 rp (13.8)
s
The fact the value of 1.44 can be raised to for a number of reasons explained on pages 405406
in the book. The following explains such a situation where we consider two satellites each with
radius rs stuck together (by gravitational forcing) so that they are separated at a distance 2rs .
The balance between tidal forcing (and not net gravity forcing as in the book) and binding now
becomes:  1/3
2p s p
3
2rs = 2
d = 2.52 rp (13.9)
d (2rs ) s
and this answer is about right, that is, if you include oblateness and rotation for the satellite
in the problem then the correct answer (d = 2.456 etc) is found. But even this situation is an
assumption because real moons will resist destruction by tidal forcing because of their tensile
strength. Examples of Moonlets that orbit within the Roche limit of a planet are Phobos in
orbit around Mars, Metis, Adrastea and Almathea for Jupiter and Pan, Atlas, Prometheus and

170
Pandora for Saturn, Cordelia, Ophelia, Bianca and Cressida for Uranus and Naiad, Thalassa
and Despina for Neptune. Over time these moonlets will disappear because the most likely
scenario is that they lose altitude so that the tidal forcing will increase.
Other examples of objects that are destroyed due to tidal forcing are comets. Shoemaker
Levy 9 approached Jupiter within the Roche limit and several other comets have been torn apart
by the tidal field of the Sun.

13.3 Exercises
The planetary sciences book [11] has various problems rated to the Hill sphere and the Roche
limit,

1. Show that our moon is in our Hill sphere

2. How long will it take before our moon reaches the Hill sphere radius with the current rate
of recession of 3 cm per jaar observed by lunar laser ranging.

3. Attempt to estimate the density ratio of a moonlet with the help of the Roche limit

171
Chapter 14

Tide generating force

The variation in gravitational pull exerted on the Earth by the motion of Sun and Moon and the
rotation of the Earth is responsible for long waves in the Earths ocean which we call tides.
On most places on Earth we experienced tides as a twice daily phenomenon where water levels
vary between a couple of decimeters to a few meters. In some bays a funneling effect takes place,
and water levels change up to 10 meter. Tides are the longest waves known in oceanography;
due to their periodicity they can be predicted well ahead in time. Tides will not only play a role
in modeling the periodic rise and fall of sea level caused by lunar and solar forcing. There are
also other phenomena that are directly related to the forcing by Sun and Moon.

14.1 Introduction
It was Newtons Principia (1687) suggesting that the difference between the gravitational at-
traction of the Moon (and the Sun) on the Earth and the Earths center are responsible for
tides, see also figure 14.1. According to this definition of astronomical tides the corresponding
acceleration f becomes:
f = f P M f EM (14.1)
whereby f P M and f EM are caused by the gravitational attraction of the Moon M. Imple-
mentation of eq. (14.1) is as straightforward as computing the lunar ephemeris and evaluating
Newtons gravitational law. In practical computations this equation is not applied because it is
more convenient to involve a tide generating potential U whose gradient U corresponds to f
in eq. (14.1).

14.2 Tide generating potential

To derive U a we start with a Taylor series of U = M /r developed at point E in figure 14.1
where M is the Moons gravitational constant and r the radius of a vector originating at point
M . The first-order approximation of this Taylor series is:

2 0 0 x1
M
f = 3 0 1 0 x2 (14.2)
rEM
0 0 1 x3

172
 fEM P
fPM
rPM
rE f

E f rEM
EM M

Figure 14.1: The external gravitational force is separated in two components, namely f EM and
f P M whose difference is according to Newtons principia (1687) responsible for the tidal force
f . Knowledge of the Earths radius rE , the Earth-Moon distance rEM and the angle is
required to compute a tide generating potential U a whose gradient U a corresponds to a tidal
acceleration vector f .

where the vector (x1 , x2 , x3 )T is originating at point E and whereby x1 is running from E
to M. The proof of equation (14.2) is explained in the following.

14.2.1 Proof

Let

U=
r
and
r = (x21 + x22 + x23 )1/2
We find that:
U
= 3 xi , i = 1, , 3
xi r
and that:
2U
= 3 5 xi xj ij 3
xi xj r r
where ij is the Kronecker symbol. Here U a originates from point M and we obtain f by
linearizing at:
x1 = r, x2 = x3 = 0
so that:
2 0 0
2U


 = 3 0 1 0
xi xj x=(r,0,0)T r
0 0 1

173
A first-order approximation of f is U |(r,0,0)T at x1 = r, x2 = x3 = 0:

2
 2 0 0 x1
U  
U |(r,0,0)T = xj = 3 0 1 0 x2
xi xj (r,0,0)T
 r
0 0 1 x3
where xi for i = 1, , 3 are small displacements at the linearization point E.

14.2.2 Work integral

We continue with equation (14.2) to derive the tide generating potential U a by evaluation of the
work integral: Z rE
Ua = (f , n) ds (14.3)
s=0
under the assumption that U a is evaluated on a sphere with radius rE .

Why a work integral?

A work integral like in eq (14.3) obtains the required amount of Joules to move from A to
B through a vector field. An example is cycling against the wind which often happens in
the Dutch climate. The cyclist goes along a certain path and n is the local unit vector in an
arbitrary coordinate system. The wind exerts a force f , and when each infinitesimal part ds is
multiplied by the projection of the wind force on n we obtain the required (or provided) work by
the wind. For potential problems we deal with a similar situation, except that the force must be
replaced by its mass-free equivalent called acceleration and where the acceleration is caused by
a gravity effect. In this case the outcome of the work integral yields potential energy difference
per mass, which is referred to as potential difference.

Evaluating the work integral

In our case n dictates the direction. Keeping in mind the situation depicted in figure 14.1 a
logical choice is:
cos
n = sin (14.4)
0
and
x1 s cos
x2 = s sin (14.5)
x3 0
so that (f , n) becomes:

2s cos cos
M
(f , n) = 3 s sin . sin
rEM
0 0
sM
2 cos2 sin2

= 3
rEM
sM 
3 cos2 1

= 3
rEM

174
It follows that:
Z rE
a sM  2

U = 3 3 cos 1 .ds
s=0 rEM
2 
M rE

3 2 1
= 3 cos (14.6)
rEM 2 2
2
M rE
= 3 P2 (cos )
rEM
which is the first term in the Taylor series where P2 (cos ) is the Legendre function of degree
2. More details on the definition of these special functions are provided in chapter 3. But there
are more terms, essentially because eq. (14.6) is of first-order. Another example is:
3U xj xk
fi = (14.7)
xi xj xk 3!
where U = /r for i, j, k = 1, , 3. Without further proof we mention that the second term in
the series derived from eq. (14.7) becomes:
3
M rE
a
Un=3 = 4 P3 (cos ) (14.8)
rEM
By induction one can show that:

M X rE n

a
U = Pn (cos ) (14.9)
rEM rEM
n=2

represents the full series describing the tide generating potential U a . In case of the Earth-Moon
1
system rE 60 rEM so that rapid convergence of eq. (14.9) is ensured. In practice it doesnt
make sense to continue the summation in eq. (14.9) beyond n = 3.

Equilibrium tides
Theoretically seen eq. (14.9) can be used to compute tidal heights at the surface of the Earth. In
a simplified case one could compute the tidal height as = g 1 U a where g is the acceleration
of the Earths gravity field. Also this statement is nothing more than to evaluate the work
integral Z Z
(f , n) ds = g ds = g = U a
0 0
assuming that g is constant. Tides predicted in this way are called equilibrium tides, they
are usually associated with Bernoilli rather than Newton who published the subject in the
Philosophae Naturalis Principea Mathematica, see also [7]. The equilibrium tide theory assumes
that ocean tides propagates with the same speed as celestrial bodies move relative to the Earth.
In reality this is not the case,
later we will show that the ocean tide propagate at a speed that
can be approximated by g.H where g is the gravitational acceleration and H the local depth
of the ocean. It turns out that our oceans are not deep enough to allow diurnal and semi-diurnal
tides to remain in equilibrium. Imagine a diurnal wave at the equator, its wavespeed would be
equal to 40 106 /(24 3600) = 463 m/s. This corresponds to an ocean with a depth of 21.5
km which exceeds an average depth of about 3 to 5 km so that equilibrium tides dont occur.

175
14.2.3 Example
In the following example we will compute g 1 (M /rEM ) (rE /rEM )n , ie. the maximum vertical
displacement caused by the tide generating potential caused by Sun and Moon. Reference values
used in equation (14.9) are (S:Sun, M:Moon):

M 4.90 1012 m3 s2 rEM 60 rE

S 1.33 1020 m3 s2 rES 1.5 1011 m
rE 6.40 106 m g 9.81 ms2

The results are shown in table 14.1.

n=2 n=3
Moon 36.2 0.603
Sun 16.5 0.703 103

Table 14.1: Displacements caused by the tide generating potential of Sun and Moon, all values
are shown in centimeters.

14.2.4 Some remarks

At the moment we can draw the following conclusions from eq. (14.9):

The P2 (cos ) term in the equation (14.9) resembles an ellipsoid with its main bulge
pointing towards the astronomical body causing the tide. This is the main tidal effect which
is, if caused by the Moon, at least 60 times larger than the n = 3 term in equation (14.9).

Sun and Moon are the largest contributors, tidal effects of other bodies in the solar system
can be ignored.

U a is unrelated to the Earths gravity field. Also it is unrelated to the acceleration expe-
rienced by the Earth revolving around the Sun. Unfortunately there exist many confusing
popular science explanations on this subject.

The result of equation (14.9) is that astronomical tides seem to occur at a rate of 2 highs
and 2 lows per day. The reason is of course Earth rotation since the Moon and Sun only
move by respectively 13 and 1 per day compared to the 359.02 per day caused by
the Earths spin rate.

Astronomical tides are too simple to explain what is really going on in nature, more on
this issue will be explained other chapters.

14.3 Frequency analysis of observed tides

Since equation (14.9) mainly depends on the astronomical positions of Sun and Moon it is not
really suitable for applications where the tidal potential is required. A more practical approach
was developed by Darwin (1883), for references see [7], who invented the harmonic method

176
of tidal analysis and prediction. It should be noted that Darwins harmonic method closely
resembles the frequency analysis method of the French mathematician and physicist Joseph
Fourier (1768-1830). Fouriers method has a general application in science and technology. The
implementations of Darwin and Doodson are dedicated to tides. Fouriers method is rather
general and can be found in several text book, a summary of the main elements of the method
can be found in appendix 4.

14.3.1 Darwin symbols and Doodson numbers

Darwins classification scheme assigns letter-digit combinations, also known as Darwin sym-
bols, to certain main lines in a spectrum of tidal lines. The M2 symbol is a typical example;
it symbolizes the most energetic tide caused by the Moon at a twice daily frequency. Later in
1921, Doodson calculated an extensive table of spectral lines which can be linked to the original
Darwin symbols. With the advent of computers in the seventies, Cartwright and Edden (1973),
with a reference to Cartwright and Tayler (1971) (hereafter CTE) for certain details, computed
new tables to verify the earlier work of Doodson. (More detailed references can be found in [6]
and in [7]). The tidal lines in these tables are identified by means of so-called Doodson numbers
D which are computed in the following way:

D = k1 (5 + k2 )(5 + k3 ).(5 + k4 )(5 + k5 )(5 + k6 ) (14.10)

where each k1 , ..., k6 is an array of small integers, corresponding with the description shown in
table 14.2, where 50 s are added to obtain a positive number. For ki = 5 where i > 0 one uses an
X and for ki = 6 where i > 0 one uses an E. In principle there exist infinitely many Doodson
numbers although in practice only a few hundred lines remain. To simplify the discussion we
divide the table in several parts: a) All tidal lines with equal k1 , which is the same as the order
m in spherical harmonics, are said to form species. Tidal species indicated with m = 0, 1, 2
correspond respectively to long period, daily and twice-daily effects, b) All tidal lines with equal
k1 and k2 terms are said to form groups, c) And finally all lines with equal k1 , k2 and k3 terms
are said to form constituents. In reality it is not necessary to go any further than the constituent
level so that a year worth of tide gauge data can be used to define amplitude and phase of a
constituent. In order to properly define the amplitude and phase of a constituent we need to
define nodal modulation factors which will be explained in chapter 17.

14.3.2 Tidal harmonic coefficients

An example of a table with tidal harmonics is shown in section 14.4. Tables 14.3 and 14.4 contain
tidal harmonic coefficients computed under the assumption that accurate planetary ephemeris
are available. In reality these planetary ephemeris are provided in the form Chebyshev polyno-
mial coefficients contained in the files provided by for instance the Jet Propulsion Laboratory
in Pasadena California USA.
To obtain the tidal harmonics we rely on a method whereby the Doodson numbers are
prescribed rather than that they are selected by filtering techniques as in CTE. We recall that
the tide generating potential U can be written in the following form:

re n
 
a M X
U = Pn (cos ) (14.11)
rem rem
n=2,3

177
The first step in realizing the conversion of equation (14.11) is to apply the addition theorem
on the Pn (cos ) functions which results in the following formulation:
n X
1
X X m (re /rem )n
Ua = Y nma (m , m )Y nma (p , p ) (14.12)
(2n + 1)rem
n=2,3 m=0 a=0

For details see chapter 3. Eq. (14.12) should now be related to the CTE equation for the tide
generating potential:
X3 X n
Ua = g cnm (p , t)fnm Pnm (cos p ) (14.13)
n=2 m=0

where g = /Re2 and for (n + m) even:

X
cnm (p , t) = H (v) [cos(Xv ) cos(mp ) sin(Xv ) sin(mp )] (14.14)
v

while for (n + m) odd:

X
cnm (p , t) = H (v) [sin(Xv ) cos(mp ) + cos(Xv ) sin(mp )] (14.15)
v

where it is assumed that:

fnm = (2Nnm )1/2 (1)m (14.16)
and:
2 (n + m)!
Nnm = (14.17)
(2n + 1) (n m)!
whereby it should be remarked that this normalization operator differs from the one used in
chapter 3. We must also specify the summation over the variable v and the corresponding
definition of Xv . In total there are approximately 400 to 500 different terms in the summation
of v each consisting of a linear combination of six astronomical elements:

Xv = k1 w1 + k2 w2 + k3 w3 + k4 w4 k5 w5 + k6 w6 (14.18)

where k1 . . . k6 are integers and:

w2 = 218.3164 + 13.17639648 T
w3 = 280.4661 + 0.98564736 T
w4 = 83.3535 + 0.11140353 T
w5 = 125.0445 - 0.05295377 T
w6 = 282.9384 + 0.00004710 T
where T is provided in Julian days relative to January 1, 2000, 12:00 ephemeris time. (When
working in UT this reference modified Julian date equals to 51544.4993.) Finally w1 is computed
as follows:
w1 = 360 U + w3 w2 180.0
where U is given in fractions of days relative to midnight. In tidal literature one usually finds the
classification of w1 to w6 as is shown in table 14.2 where it must be remarked that w5 is retrograde
whereas all other elements are prograde. This explains the minus sign equation (14.18).

178
 Here Frequency Cartwright, Explanation
Doodson
k1 ,w1 daily , mean time angle in lunar days
k2 ,w2 monthly q, s mean longitude of the moon
k3 ,w3 annual q0, h mean longitude of the sun
k4 ,w4 8.85 yr p, p mean longitude of lunar perigee
k5 ,w5 18.61 yr N , N 0 mean longitude of ascending lunar node
k6 ,w6 20926 yr p0 , p1 mean longitude of the sun at perihelion

Table 14.2: Classification of frequencies in tables of tidal harmonics. The columns contain: [1]
the notation used in the Doodson number, [2] the frequency, [3] notation used in tidal literature,
[4] explanation of variables.

14.4 Tidal harmonics

Section 14.3.2 introduced the concept of tidal harmonics. Purpose of this section is to present the
implementation of a method to obtain the tables and to present the results. The method used
here to compute tidal harmonics in Cartwright Tayler and Edden differs from the approach used
in this lecture notes. In contrast to CTE, who used several convolution operators to separate
tidal groups. Here we rely on an algorithm that assumes a least squares fitting procedure and
prior knowledge of all Doodson numbers in the summation over all frequencies indicated by
index v. To obtain the tidal harmonic coefficients H (v) for each Doodson number the following
procedure is used:
For each degree n and tidal species m (which equals k1 ) the algorithm starts to collect all
matching Doodson numbers.

The following step is to generate values of:

a b (re /reb (t))n

Unm (t) = Pnm (cos b (t)) cos(mb (t))
(2n + 1)reb (t)
where t is running between 1990/1/1 00:00 and 2010/1/1 00:00 in a sufficiently dense
number of steps to avoid under sampling. Positions of Sun and Moon obtained from a
planetary ephemeris model are used to compute the distance Reb (t) between the astro-
nomical body (indicated by subscript b) and the Earths center (indicated by subscript e)
are transformed into Earth-fixed coordinates to obtain b (t) and b (t).

The following step is a least squares analysis of Unm (t) where the observations equations
are as follows:
0
X
a
Unm (t) = G(v ) cos(Xv0 )
v0
when m + n is even and
0
X
a
Unm (t) = G(v ) sin(Xv0 )
v0

v0
whenever m + n is odd. The symbol is used to indicate that we are only considering the
appropriate subset of Doodson numbers to generate the Xv0 values, see also section 14.3.2.

179
 0
Finally the Gv values need a scaling factor to convert them into numbers that have
the same dimension as one finds in CTE. Partly this conversion is caused by a different
normalization between surface harmonics used in CTE and eqns. (14.13), (14.14) and
(14.15) here, although is it also required to take into account the factor g. As a result:
0 0
H v = Gv g 1 fnm
1 2
nm
where nm is the normalization factor as used in chapter 3 and fnm the normalization
factor used by CTE given in eqns. (14.16) and (14.17). In our algorithm g is computed as
/re2 where = 3.9860044 1014 [m3 /s2 ] and re = 6378137.0 [m].
For all collected spectral lines we show in table 14.3 and 14.4 only those where |H (v) | exceeds the
value of 0.0025. Tables 14.3 and 14.4 show in columns 2 to 7 the values of k1 till k6 , in column
8 the degree n, in column 9 the coefficient H v in equations (14.14) and (14.15), in column 10
the Darwin symbol provided that it exists, and in column 11 the Doodson number.
Some remarks about the tables: a) The tables only hold in the time period indicated earlier
in this chapter, b) There are small differences, mostly in the 5th digit behind the period, with
respect to the values given in [6], c) In total we have used 484 spectral lines although many more
tidal lines may be observed with a cryogenic gravimeter.

14.5 Exercises
1. Show that the potential energy difference for 0 to H meter above the ground becomes
m.g.H kg.m2 /s2 . Your answer must start with the potential function U = /r.
2. Show that the outcome of Newtons gravity law for two masses m1 and m2 evaluated for
one of the masses corresponds to the gradient of a so-called point mass potential function
U = G.m1 /r + const. Verify that the point mass potential function in 3D exactly fullfills
the Laplace equation.
3. Show that the function 1/rP M in figure 14.1 can be developed in a series of Legendre
functions Pn (cos ).
4. Show that a work integral for a closed path becomes zero when the force is equal to a mass
times an acceleration for a potential functions that satisfy the Laplace equation.
5. Show that a homogeneous hollow sphere and a solid equivalent generate the same potential
field outside the sphere.
6. Compute the ratio between the acceleration terms Fem and Fpm in figure 14.1 at the
Earths surface. Do this at the Poles and the Lunar sub-point. Example 14.2.3 provides
constants that apply to the Earth Moon Sun problem.
7. Assume that the astronomical tide generating potential is developed to degree 2, for which
values of is the equilibrium tide zero?
8. Compute the extreme tidal height displacements for the equilibrium tide on Earth caused
by Jupiter, its mass ratio with respect to Earth is 317.8.
9. How much observation time is required to separate the S2 tide from the K2 tide.

180
 k1 k2 k3 k4 k5 k6 n H (v) Darwin Doodson
1 0 0 0 0 0 0 2 -.31459 M0 + S 0 055.555
2 0 0 0 0 1 0 2 .02793 055.565
3 0 0 1 0 0 -1 2 -.00492 Sa 056.554
4 0 0 2 0 0 0 2 -.03099 Ssa 057.555
5 0 1 -2 1 0 0 2 -.00673 063.655
6 0 1 0 -1 -1 0 2 .00231 065.445
7 0 1 0 -1 0 0 2 -.03518 Mm 065.455
8 0 1 0 -1 1 0 2 .00228 065.465
9 0 2 -2 0 0 0 2 -.00584 073.555
10 0 2 0 -2 0 0 2 -.00288 075.355
11 0 2 0 0 0 0 2 -.06660 Mf 075.555
12 0 2 0 0 1 0 2 -.02761 075.565
13 0 2 0 0 2 0 2 -.00258 075.575
14 0 3 -2 1 0 0 2 -.00242 083.655
15 0 3 0 -1 0 0 2 -.01275 085.455
16 0 3 0 -1 1 0 2 -.00529 085.465
17 0 4 -2 0 0 0 2 -.00204 093.555
18 1 -3 0 2 0 0 2 .00664 125.755
19 1 -3 2 0 0 0 2 .00801 1 127.555
20 1 -2 0 1 -1 0 2 .00947 135.645
21 1 -2 0 1 0 0 2 .05019 Q1 135.655
22 1 -2 2 -1 0 0 2 .00953 1 137.455
23 1 -1 0 0 -1 0 2 .04946 145.545
24 1 -1 0 0 0 0 2 .26216 O1 145.555
25 1 -1 2 0 0 0 2 -.00343 147.555
26 1 0 0 -1 0 0 2 -.00741 155.455
27 1 0 0 1 0 0 2 -.02062 M1 155.655
28 1 0 0 1 1 0 2 -.00414 155.665
29 1 0 2 -1 0 0 2 -.00394 157.455
30 1 1 -3 0 0 1 2 .00713 1 162.556
31 1 1 -2 0 0 0 2 .12199 P1 163.555
32 1 1 -1 0 0 1 2 -.00288 S1 164.556
33 1 1 0 0 -1 0 2 .00730 165.545
34 1 1 0 0 0 0 2 -.36872 K1 165.555

Table 14.3: Tidal harmonic constants

181
 k1 k2 k3 k4 k5 k6 n H (v) Darwin Doodson
35 1 1 0 0 1 0 2 -.05002 165.565
36 1 1 1 0 0 -1 2 -.00292 1 166.554
37 1 1 2 0 0 0 2 -.00525 1 167.555
38 1 2 -2 1 0 0 2 -.00394 1 173.655
39 1 2 0 -1 0 0 2 -.02062 J1 175.455
40 1 2 0 -1 1 0 2 -.00409 175.465
41 1 3 -2 0 0 0 2 -.00342 183.555
42 1 3 0 0 0 0 2 -.01128 OO1 185.555
43 1 3 0 0 1 0 2 -.00723 185.565
44 1 4 0 -1 0 0 2 -.00216 195.455
45 2 -3 2 1 0 0 2 .00467 227.655
46 2 -2 0 2 0 0 2 .01601 2N2 235.755
47 2 -2 2 0 0 0 2 .01932 2 237.555
48 2 -1 0 1 -1 0 2 -.00451 245.645
49 2 -1 0 1 0 0 2 .12099 N2 245.655
50 2 -1 2 -1 0 0 2 .02298 2 247.455
51 2 0 -1 0 0 1 2 -.00217 254.556
52 2 0 0 0 -1 0 2 -.02358 255.545
53 2 0 0 0 0 0 2 .63194 M2 255.555
54 2 1 -2 1 0 0 2 -.00466 263.655
55 2 1 0 -1 0 0 2 -.01786 L2 265.455
56 2 1 0 1 0 0 2 .00447 265.655
57 2 2 -3 0 0 1 2 .01719 T2 272.556
58 2 2 -2 0 0 0 2 .29401 S2 273.555
59 2 2 -1 0 0 -1 2 -.00246 274.554
60 2 2 0 0 0 0 2 .07992 K2 275.555
61 2 2 0 0 1 0 2 .02382 275.565
62 2 2 0 0 2 0 2 .00259 275.575
63 2 3 0 -1 0 0 2 .00447 285.455
64 0 1 0 0 0 0 3 -.00375 065.555
65 1 0 0 0 0 0 3 .00399 155.555
66 2 -1 0 0 0 0 3 -.00389 245.555
67 2 1 0 0 0 0 3 .00359 265.555
68 3 -1 0 1 0 0 3 -.00210 345.655
69 3 0 0 0 0 0 3 -.00765 355.555

Table 14.4: Tidal harmonic constants

182
Chapter 15

Tides deforming the Earth

Imagine that the solid Earth itself is somehow deforming under tidal accelerations, i.e. gradients
of the tide generating potential. This is not unique to our planet, all bodies in the universe
experience the same effect. Notorious are moons in the neighborhood of the larger planets such
as Saturn where the tidal forces can exceed the maximum allowed stress causing the Moon to
collapse.
It must be remarked that the Earth will resist forces caused by the tide generating potential.
This was recognized by A.E.H. Love (1927), see [6], who assumed that an applied astronomical
tide potential for one tidal line:
X X
Ua = Una = Un0 (r)Sn e( jt) (15.1)
n n

where Sn is a surface harmonic, will result in a deformation at the surface of the Earth:
un (R) = g 1 [hn (R)Sn er + ln (R)Sn et ] Un0 (R)e( jt) (15.2)
where er and et are radial and tangential unit vectors. The indirect potential caused by this
solid Earth tide effect will be:
U (R) = kn (R)Un0 (R)Sn e( jt) (15.3)
Equations (15.2) and (15.3) contain so-called Love numbers hn , kn and ln describing the geo-
metric radial, indirect potential and geometric tangential effects. Finally we remark that
Love numbers can be obtained from geophysical Earth models and also from geodetic space tech-
nique such as VLBI, see table 15.1 taken from [31], where we present the Love numbers reserved
for the deformations by a volume force, or potential, that does not load the surface. Loading is
described by separate Love numbers h0n , kn0 and ln0 that will be discussed in chapter 18.

15.1 Solid Earth tides

According to equations (15.2) and (15.3) the solid Earth itself will deform under the tidal
forces. Well observable is the vertical effect resulting in height variations at geodetic stations.
To compute the so-called solid-Earth tide s we represent the tide generating potential as the
series:
X
Ua = Una
n=2

183
 Dziewonski-Anderson Gutenberg-Bullen
n hn kn ln hn kn ln
2 0.612 0.303 0.0855 0.611 0.304 0.0832
3 0.293 0.0937 0.0152 0.289 0.0942 0.0145
4 0.179 0.0423 0.0106 0.175 0.0429 0.0103

Table 15.1: Love numbers derived from the Dziewonski-Anderson and the Gutenberg-Bullen
Earth models.

length NS baselines EW baselines

1 0.003 0.004
2 0.006 0.009
5 0.016 0.022
10 0.031 0.043
20 0.063 0.084
50 0.145 0.186
90 0.134 0.237

Table 15.2: The maximum solid earth tide effect [m] on the relative vertical coordinates of
geodetic stations for North-South and East-West baselines varying in length between 0 and 90
angular distance.

so that:

X
1
s = g hn Una (15.4)
n=2

An example of s is shown in table 15.2 where the extreme values of |s | are tabulated as a
relative height of two geodetic stations separated by a certain spherical distance. One may
conclude that regional GPS networks up to e.g. 200 by 200 kilometers are not significantly
affected by solid earth tides; larger networks are affected and a correction must be made for the
solid Earth tide. The correction itself is probably accurate to within 1 percent or better so that
one doesnt need to worry about errors in excess of a couple of millimeters.

15.2 Long period equilibrium tides in the ocean

At periods substantially longer than 1 day the oceans are in equilibrium with respect to the
tide generating potential. But also here the situation is more complicated than one immediately
expects from equation (14.9) due to the existence of kn in equation (15.3). For this reason long
period equilibrium tides in the oceans are derived by:
X
e = g 1 (1 + kn hn )Una (15.5)
n

184
essentially because the term (1 + kn ) dictates the geometrical shape of the oceans due to the
tide generating potential but also the indirect or induced potential kn Una . Still there is a need
to include hn Una since ocean tides are always relative to the sea floor or land which is already
experiencing the solid earth tide effect s described in equation (15.4). Again we emphasize that
equation (15.5) is only representative for a long periodic response of the ocean tide which is in
a state of equilibrium. Hence equation (15.5) must only be applied to all m = 0 terms in the
tide generating potential.

15.3 Tidal accelerations at satellite altitude

The astronomical tide generating potential U at the surface of the Earth with radius re has the
usual form:

p X p X
U (re ) = (re /rp )n Pn (cos ) = (re /rp )n+1 Pn (cos ) (15.6)
rp re
n=2 n=2

The potential can also be used directly at the altitude of the satellite to compute gradients, but
in fact there is no need to do this since the accelerations can be derived from Newtons definition
of tidal forces. This procedure does not anymore work for the induced or secondary potential
U 0 (re ) since the theory of Love predicts that:

p X
U 0 (re ) = (re /rp )n+1 kn Pn (cos ) (15.7)
re
n=2

where it should be remarked that this expression is the result of a deformation of the Earth as a
result of tidal forcing. The effect at satellite altitude should be that of an upward continuation,
in fact, it is a mistake to replace re by the satellite radius rs in the last equation. Instead to
bring U 0 (re ) to U 0 (rs ) we get the expression:

p X
0
U (rs ) = (re /rs )n+1 (re /rp )n+1 kn Pn (cos ) (15.8)
re
n=2

Finally we eliminate cos() by use of the addition theorem of Legendre functions:

 n+1 n
0 p X re2 kn X
U (rs ) = P nm (cos p )P nm (cos s ) cos(m(s p )) (15.9)
re rs rp 2n + 1
n=2 m=0

where (rs , s , s ) and (rp , p , p ) are spherical coordinates in the terrestial frame respectively
for the satellite and the planet in question. This is the usual expression as it can be found in
literature, see for instance [31].
Gradients required for the precision orbit determination (POD) software packages are derived
from U (rs ) and U 0 (rs ) first in spherical terrestial coordinates which are then transformed via
the appropriate Jacobians into terrestial Cartesian coordinates and later in inertial Cartesian
coordinates which appear in the equations of motion in POD. Differentiation rules show that
the latter transformation sequence follows the transposed transformation sequence compared to
that of vectors.

185
 Satellite orbit determination techniques allow one to obtain in an indepent way the k2 Love
number of the Earth or of an arbitrary body in the solar system. Later in these notes it will
be shown that similar techniques also allow to estimate the global rate of dissipation of tidal
energy, essentially because tidal energy dissipation result in a phase lag between the tidal bulge
and the line connecting the Earth to the external planet for which the indirect tide effect is
computed.

15.4 Gravimetric solid earth tides

A gravimeter is an instrument for observing the actual value of gravity. There are several types of
instruments, one type measures gravity difference between two locations, another type measures
the absolute value of gravity. The measured quantity is usually expressed in milligals (mgals)
relative to an Earth reference gravity model. The milligal is not a S.I. preferred unit, but it is
still used in research dealing with gravity values on the Earths surface, one mgal equals 105
m/s2 , and the static variations referring to a value at the mean sea level vary between -300 to
+300 mgal. Responsible for these static variations are density anomalies inside the Earth.
Gravimeters do also observe tides, the range is approximately 0.1 of a mgal which is within
the accuracy of modern instruments. Observed are the direct astronomical tide, the indirect
solid earth tide but also the height variations caused by the solid Earth tides. According to [38]
we have the following situation:
V0
V = V0 + s + Ua + UI (15.10)
r
where V is the observed potential, V0 is the result of the Earths gravity field, s the vertical
displacement implied by the solid Earth tide, U a is the tide generating potential and U i the
indirect solid Earth tide potential. In the following we assume that:
X  r n
Ua = Una
n
r0

X  r0 n+1
Ui = kn Una
n
r
V
= 2 = g
r r
where is the Earths gravitational constant, r0 the mean equatorial radius, and Una the tide
generating potential at r0 . Note that in the definition of the latter equation we have taken the
potential as a negative function on the Earth surface where attains a positive value. This is
also the correct convention since the potential energy of a particle must be increased to lift it
from the Earth surface and it must become zero at infinity. We get:

V V0 2V U a U i
= + s 2 + +
r r r r r
which becomes:
 n
V V0 2g X n r X (n + 1)  r0 n+1
= + s + Una kn Una
r r r n
r r0 n
r r

186
where 2 V /r2 is approximated by 2g/r assuming a point mass potential function. When
substituting the solid Earth tide effect s we get:

V V0 2g X X  n   r n X (n + 1)  r0 n+1
= + hn Una g 1 + Una kn Una
r r r n n
r r0 n
r r

so that for r r0 :
nUna
   
V V0 X 2hn n+1
= + +1 kn
r r n
n n r
which becomes:
X Una
  
2 n+1
g = g0 + 1 + hn kn
n
n n r
On gravity anomalies the effect becomes:
X Una
  
2 n+1
g = g g0 = 1 + hn kn
n
n n r

The main contribution comes from the term:

U2a U a
 
3
g = 1 + h2 k2 = 1.17 2
2 r r

while a secondary contribution comes from the term:

U3a U a
 
2 4
g = 1 + h3 k3 = 1.07 3
3 3 r r

This shows that gravimeters in principle sense a scaled version of the astronomic tide potential,
the factors 1.17 and 1.07 are called gravimetric factors. By doing so gravimetric observations
add their own constraint to the definition of the Love numbers h2 and k2 and also h3 and k3 .

15.5 Reference system issues

In view of equation (15.5) we must be careful in defining parameters modeling the reference
ellipsoid. The reason is due to a contribution of the tide generating potential at Doodson
number 055.555 where it turns out that:

g 1 U2a = 0.19844 P2,0 (sin ) (15.11)

1
g k2 U2a = 0.06013 P2,0 (sin ) (15.12)
1
g (1 + k2 )U2a = 0.25857 P2,0 (sin ) (15.13)

where we have assumed that k2 = 0.303, h2 = 0.612 and H (v) = 0.31459 at Doodson number
055.555. The question which equation goes where is not as trivial as one might think. In
principle there are three tidal systems, and the definition is as follows:

A tide free system: this means that eqn. (15.13) is removed from the reference ellipsoid
flattening.

187
 A zero-tide system: this means that eqn. (15.11) is removed but that (15.12) is not removed
from the reference ellipsoid flattening.

A mean-tide system: this means that eqns. (15.13) is not removed from the reference
ellipsoid.

Important in the discussion is that the user of a reference system must be aware which choice
has been made in the definition of the flattening parameter of the reference ellipsoid. The
International Association of Geodesy recommends a zero-tide system so that it is not necessary
to define k2 at the zero frequency. In fact, from a rheologic perspective it is unclear which value
should be assigned to k2 , the IAG recommendation is therefore the most logical choice.

15.6 Exercises
1. Show that the Love numbers h2 and k2 can be estimated from observations of the gravime-
ter tide in combination with observations of the long periodic ocean tide observed by tide
gauges.

2. What are the extreme variations in the water level of the M2 equilibrium tide at a latitude
of 10N.

3. What are the extreme variations in mgal of the M2 gravimetric tide at a latitude of 50S.

4. What is the largest relative gravimetric tidal effect between Amsterdam and Paris as a
result of the Moon.

5. Verify equation (15.11), how big is this effect between Groningen and Brussel.

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Chapter 16

Ocean tides

Purpose of this chapter is to introduce some basic properties concerning the dynamics of fluids
that is applicable to the ocean tide problem. Of course the oceans themselves will respond
differently to the tide generating forces. Ocean tides are exactly the effect that one observes
at the coast; i.e. the long periodic, diurnal and semi-diurnal motions between the sea surface
and the land. In most regions on Earth the ocean tide effect is approximately 0.5 to 1 meters
whereas in some bays found along the coast of e.g. Normandy and Brittany the tidal wave is
amplified to 10 meters. Ocean tides may have great consequences for daily life and also marine
biology in coastal areas. Some islands such as Mt. Saint Michele in Brittany cant be reached
during high tide if no separate access road would exist. A map of the global M2 ocean tide
is given in figure 16.1 from which one can see that there are regions without any tide which
are called amphidromes where a tidal wave is continuously rotating about a fixed geographical
location. If we ignore friction then the orientation of the rotation is determined by the balance
between the pressure gradient and the Coriolis force. It was Laplace who laid the foundations
for modern tidal research, his main contributions were:

The separation of tides into distinct Species of long period, daily and twice daily (and
higher) frequencies.

The (almost exact) dynamic equations linking the horizontal and vertical displacement of
water particles with the horizontal components of the tide-raising force.

The hypothesis that, owing to the dominant linearity of these equations, the tide at any
place will have the same spectral frequencies as those present in the generating force.

Laplace derived solutions for the dynamic equations only for the ocean and atmospheres covering
a globe, but found them to be strongly dependent on the assumed depth of fluid. Realistic
bathymetry and continental boundaries rendered Laplaces solution mathematically intractable.
To explain this problem we will deal with the following topics:

Define the equations of motion

What is advection, friction and turbulence

The Navier Stokes equations

Laplace tidal equations

189
Figure 16.1: The top panel shows the amplitudes in centimeter of the M2 ocean tide, the bottom
panel shows the corresponding phase map.

190
 A general wave solution, the Helmholtz equation
Dispersion relations
However we will avoid to represent a complete course in physical oceanography; within the scope
of this course on tides we have to constrain ourselves to a number of essential assumptions and
definitions.

16.1 Equations of motion

16.1.1 Newtons law on a rotating sphere
The oceans can be seen as a thin rotating shell with a thickness of approximately 5 km relative
on a sphere with an average radius of 6371 km. To understand the dynamics of fluids in this thin
rotating shell we initially consider Newtons law f = m.a for a given water parcel at a position:

x = ei xi = ea xa (16.1)

In this equation ei and ea are base vectors. Here the i index is used for the inertial coordinate
frame, the local Earth-fixed coordinate system gets index a. Purpose of the following two sections
will be to find expressions for inertial velocities and accelerations and their expressions in the
Earth fixed system, which will appear in the equations of motion in fluid dynamics.

Inertial velocities and accelerations

There is a unique relation between the inertial and the Earth-fixed system given by the trans-
formation:
ei = Ria ea (16.2)
In the inertial coordinate system, velocities can be derived by a straightforward differentiation
so that:
x = ei xi (16.3)
and accelerations are obtained by a second differentiation:

x = ei xi (16.4)

Note that this approach is only possible in an inertial frame, which is a frame that does not rotate
or accelerate by itself. If the frame would accelerate or rotate then ei also contains derivatives
with respect to time. This aspect is worked out in the following section.

Local Earth fixed velocities and accelerations

The Earth fixed system is not an inertial system due to Earth rotation. In this case the base
vectors themselves follow different differentiation rules:

ea = ea (16.5)

where denotes the vector (0, 0, ) for an Earth that is rotating about its z-axis at a constant
speed of radians per second. We find:

ea = ea + ea (16.6)

191
and:
x = ea xa + 2ea xa + ea xa (16.7)
which is equivalent to:

x = ea xa + ea xa + 2 ea xa + ea xa (16.8)

leading to the equation:

xi = xa + 2 xa + xa + xa (16.9)

where xi is the inertial acceleration vector, xa the Earth-fixed acceleration vector. The difference
between these vectors is the result of frame accelerations:

The term 2 xa is known as the Coriolis effect. Consequence of the Coriolis effect is that
particles moving over the surface of the Earth will experience an apparent force directed
perpendicular to their direction. On Earth the Coriolis force is directed to East when a
particle is moving to the North on the Northern hemisphere.

The term xa is a centrifugal contribution. This results in an acceleration component

that is directed away from the Earths spin axis.

The term xa indicates a rotational acceleration which can be ignored unless one intends
to consider the small variations in the Earths spin vector .

16.1.2 Assembly step momentum equations

To obtain the equations of motion for fluid problems we will consider all relevant accelerations
that act on a water parcel in the Earths fixed frame:

g is the sum of gravitational and centrifugal accelerations, ie. the gravity acceleration
vector,

2 u is the Coriolis effect which is an apparent acceleration term caused by Earth

rotation,

f symbolizes additional accelerations which are for instance caused by friction and advec-
tion in fluids,

1 p is the pressure gradient in a fluid.

The latter two terms are characteristic for motions of fluids and gasses on the Earths surface.
The pressure gradient is the largest, and it will be explained first because it appears in all
hydrodynamic models.

The pressure gradient

This gradient follows from the consideration of a pressure change on a parcel of water as shown
in figure 16.2. In this figure there is a pressure p acting on the western face dy.dz and a pressure
p + dp acting on the eastern face dy.dz. To obtain a force we multiply the pressure term times

192
 Figure 16.2: Pressure gradient

the area on which it is acting. The difference between the forces is only relevant since p itself
could be the result of a static situation:
p.dy.dz (p + dp)dy.dz = dpdydz
To obtain a force by volume one should divide this expression by dx.dy.dz to obtain:
p

x
To obtain a force by mass one should divide by .dx.dy.dz to obtain:
1 p

x
This expression is the acceleration of a parcel towards the East which is our x direction. To
obtain the acceleration vector of the water parcel one should compute the gradient of the pressure
field p and scale with the term 1/.

Geostrophic balance
The following expression considers the balance between local acceleration, the pressure gradient,
the Coriolis effect and residual forces f :
Du 1
= p 2 u + g + f . (16.10)
Dt
This vector equation could also be formulated as three separate equations with the local coor-
dinates x, y and z and the corresponding velocity components u, v and w. Here we follow the

193
 Figure 16.3: Choice of the local coordinate system relevant to the equations of motion.

convention found in literature and assign the x-axis direction corresponding with the u-velocity
component to the local east, the y-axis direction and corresponding v-velocity component to
the local north, and the z-axis including the w-velocity pointing out of the sea surface, see also
figure 16.3. All vectors in equation (16.10) must be expressed in the local x, y, z coordinate
frame. If corresponds to the latitude of the water parcel and to the length of then the
following substitutions are allowed:

= (0, cos , sin )T

g = (0, 0, g)T
f = (Fx , Fy , Fz )T
v = (u, v, w)T

The result after substitution is the equations of motions in three dimensions:

Du 1 p
= + Fx + 2 sin v 2 cos w
Dt x
Dv 1 p
= + Fy 2 sin u (16.11)
Dt y
Dw 1 p
= + Fz + 2 cos u g
Dt z
Providing that we forget about dissipative and advective terms eqns. (16.11) tell us nothing
more than that the pressure gradient, the Coriolis force and the gravity vector are in balance,
see also figure 16.4. Some remarks with regard to the importance of acceleration terms in
eqns. (16.11)(a-c):

194
Figure 16.4: The equations of motion is dynamical oceanography, the Coriolis force, the pressure
gradient and the gravity vector are in balance.

The vertical velocity w is small and we will drop this term.

In eq. (16.11)(c) the gravity term and the pressure gradient term dominate, cancellation
of the other terms results in the hydrostatic equation telling us that pressure linearly
increases by depth.

The term f = 2 sin is called the Coriolis parameter.

16.1.3 Advection
The terms Du/Dt, Dv/Dt and Dw/Dt in eqns. (16.11) should be seen as absolute derivatives.
In reality these expressions contain an advective contribution.

Du u u u u
= + u. + v. + w.
Dt t x y z
Dv v v v v (16.12)
= + u. + v. + w.
Dt t x y z
Dw w w w w
= + u. + v. + w.
Dt t x y z
In literature terms like u/t are normally considered as so-called local accelerations whereas
advective terms like uu/x + ... are considered as field accelerations. The physical inter-
pretation is that two types of acceleration may take place. In the first terms on the right
hand side, accelerations occur locally at the coordinates (x, y, z) resulting in u/t, v/t, and

195
w/t whereas in the second case the velocity vector is changing with respect to the coor-
dinates resulting in advection. This effect is non-linear because velocities are squared, (e.g.
u(u/x) = 21 [(u2 )/x]).

16.1.4 Friction
In eq. (16.11) friction may appear in Fx , Fy and Fz . Based upon observational evidence, Stokes
suggested that tangentional stresses are related to the velocity shear as:

ij = (ui /xj + uj /xi ) (16.13)

where is a molecular viscosity coefficient characteristic for a particular fluid. Frictional forces
are obtained by:
2 ui
 
ij ui
F = = + (16.14)
xj xj 2 xi xj
which is approximated by:
2 ui
F = (16.15)
xj 2
if an incompressible fluid is assumed. A separate issue is that viscosity = / may not be
constant because of turbulence. In this case:
 
ij ui
F = = (16.16)
xj xj xj

although it should be remarked that also this equation is based upon an assumption. As a
general rule, no known oceanic motion is controlled by molecular viscosity, since it is far too
weak. In ocean dynamics the Reynold stress involving turbulence or eddy viscosity always
applies, see also [43] or [45].

16.1.5 Turbulence
Motions of fluids often show a turbulent behavior whereby energy contained in small scale phe-
nomena transfer their energy to larger scales. In order to assess whether turbulence occurs in an
experiment we define the so-called Reynolds number Re which is a measure for the ratio between
advective and the frictional terms. The Reynolds number is approximated as Re = U.L/, where
U and L are velocities and lengths at the characteristic scales at which the motions occurs. Large
Reynolds numbers, e.g. ones which are greater than 1000, usually indicates turbulent flow.
An example of this phenomenon can be found in the Gulf stream area where L is of the order
of 100 km, U is of the order of 1 m/s and a typical value for is approximately 106 m2 s1 so
that Re = U.L/ 1011 . The effect displays itself as a meandering of the main stream which
can be nicely demonstrated by infrared images of the area showing the turbulent flow of the
Gulf stream occasionally releasing eddies that will live for considerable time in the open oceans.
The same phenomenon can be observed in other western boundary regions of the oceans such
as the Kuroshio current East of Japan and the Argulhas retroreflection current south of Cape
of Good Hope.

196
 Figure 16.5: Continuity and depth averaged velocities

16.2 Laplace Tidal Equations

So far the equations of motions are formulated in three dimensions. The goal of the Laplace Tidal
Equations is in first instance to simplify this situation. Essentially the LTE describe the motions
of a depth averaged velocity fluid dynamics problem. Rather than considering the equations of
motion for a parcel of water in three dimensions, the problem is scaled down to two dimensions
in x and y whereby the former is locally directed to the east and the latter locally directed to
the north. A new element in the discussion is a consideration of the continuity equation.
To obtain the LTE we consider a box of water with the ground plane dimensions dx times
dy and height h representing the mean depth of the ocean, see also figure 16.5. Moreover let
u1 be the mean columnar velocity of water entering the box via the dy h plane from the west
and u2 the mean velocity of water leaving the box via the dy h plane to the east. Also let v1
be the mean columnar velocity of water entering the box via the dx h plane from the south
and v2 the mean velocity of water leaving the dx h plane to the north. In case there are no
additional sources or drains (like a hole in the ocean floor or some river adding water to it) we
find that:
dV
h.dy.(u2 u1 ) + h.dx.(v2 v1 ) + =0 (16.17)
dt
where the volume V is computed as dx.dy.h. Take as the surface elevation due to the in-flux
of water and:
dV d
= dx.dy. (16.18)
dt dt

197
If the latter equation is substituted in eq.(16.17) and all terms are divided by dx.dy we find:
 
u v
h + + =0 (16.19)
x y t

The latter equation should now be combined with eq. (16.11) where the third equation can be
simplified as a hydrostatic approximation essentially telling us that a water column of meters
is responsible for a certain pressure p:
p = g.. (16.20)
following the requirement that the pressure p is computed relative to a surface that doesnt
experience a change in height. We get the horizontal pressure gradients:
1 p (g) 1 p (g)
= and = (16.21)
x x y y
Moreover for the forcing terms Fx and Fy in eq. (16.11) we substitute the horizontal gradients:

U a U a
Fx = + Gx and Fy = + Gy (16.22)
x y
where U a is the total tide generating potential and Gx and Gy terms as a result of advection
and/or friction. Substitution of eqns. (16.21) and (16.22) in eqn. (16.11) and elimination of the
term 2 cos()w in the first and second equation results in a set of equations which were first
formulated by Laplace:
Du
= (g + U a ) + f.v + Gx
Dt x
Dv
= (g + U a ) f.u + Gy (16.23)
Dt y
 
D u v
= h +
Dt x y

The Laplace tidal equations consist of two parts; equations (16.23)(a-b) are called the momentum
equations, and (16.23)(c) is called the continuity equation. Various refinements are possible, two
relevant refinements are:
We have ignored the effect of secondary tide potentials caused by ocean tides loading on
the lithosphere, more details can be found in chapter 18.

The depth term h could by replaced by h + because the ocean depth is increased by the
water level variation (although this modification would introduce a non-linearity).

For the LTE:  h.

To solve the LTE it is also necessary to pose initial and boundary conditions including a domain
in which the equations are to be solved. From physical point of view a no-flux boundary condition
is justified, in which case (u, n) = 0 with n perpendicular to the boundary of the domain. For a
global tide problem the domain is essentially the oceans, and the boundary is therefor the shore.
Other possibilities are to define a half open tide problem where a part of the boundary is
on the open ocean where water levels are prescribed while another part is closed on the shore.

198
This option is often used in civil engineering application where it is intended to study a limited
area problem. Other variants of boundary conditions including reflecting or (weakly) absorbing
boundaries are an option in some software packages.
In the next section we show simple solutions for the Laplace tidal equations demonstrating
that the depth averaged velocity problem, better known as the barotropic tide problem, can be
approximated by a Helmholtz equation which is characteristic for wave phenomena in physics.

16.3 Helmholtz equation

Intuitively we always assumed that ocean tides are periodic phenomena, but of course it would
be nicer to show under which conditions this is the case. Let us introduce a test solution for the
problem where we assume that:

u(t) = ue( jt) (16.24)

(
v(t) = ve jt) (16.25)
(
(t) = e jt) (16.26)

where j = 1. For tides we know that the gradient of the tide generating potential is:

U a (t) = e( jt) (16.27)

Furthermore we will simplify advection and friction and assume that these terms can be approx-
imated by:

Gx (t) = Gx e( jt) (16.28)

(
Gy (t) = Gy e jt) (16.29)

If this test solution is substituted in the momentum equations then we obtain:

        
j f u /x /x Gx
= g + + (16.30)
+f j v /y /y Gy

Provided that we are dealing with a regular system of equations it is possible to solve u and v
and to substitute this solution in the continuity equation that is part of the LTE. After some
manipulation we get:
! !
Gx Gy jf h Gx Gy
 2
2 2 2

2 2
( f )+gh + =h + + + 2 + (16.31)
x2 y 2 x y x2 y y x

The left hand side of equation (16.31) is known as the Helmholtz equation which is typical for
wave phenomena in physics. The term gh in eq. (16.31) contains the squared surface speed (c) of
a tidal wave. Some examples are: a tidal wave in a sea of 50 meter depth runs with a velocity of

50.g which is about 22 m/s or 81 km/h. In an ocean of 5 km depth c will rapidly increase, we
get 223.61 m/s or 805 km/h which is equal to that of an aircraft. A critical step in the derivation
of the Helmholtz equation is the treatment of advection and friction term contained in Gx and
Gy and the vorticity term . As long as these terms are written in the form of harmonic test
functions like in (16.28) and (16.29) there is no real point of concern. To understand this issue
we must address the problem of a drag law that controls the dissipation of a tidal wave.

199
16.4 Drag laws
The drag law is an essential component of a hydrodynamic tide model, omission of a dissipative
mechanism results in modeling tides as an undamped system since tidal waves can not lose their
energy. Physically seen this is completely impossible because the tides are continuously excited
by gravitational forcing. A critical step is therefor the formulation of a dissipative mechanism
which is often chosen as a bottom friction term. Friction between layers of fluid was initially
considered to be too small to explain the dissipation problem in tides, friction against the walls
of a channel or better the ocean floor is considered to be more realistic. In this way the ocean
tides dissipate more than 75 percent of their energy, more details are provided in chapter 20.
There is an empirical law for bottom drag which was found by the Frenchman Chezy who
found that drag is proportional to the velocity squared and inverse proportional to the depth
of a channel. Chezy essentially compared the height gradient of rivers against the flow in the
river and geology of the river bed. Under such conditions the river bed drag has to match the
horizontal component of the pressure gradient, which essentially follows from the height gradient
of the river. The Chezy law extended to two dimensions is:
p
Gx = Cd u u2 + v 2 (16.32)
p
Gy = Cd v u2 + v 2 (16.33)

where Cd = g/(hCz2 ), g is gravity, h is depth and Cz a scaling coefficient, or the Chezy coefficient.
In reality Cz depends on the physical properties of the river bed; reasonable values are between
40 and 70.
Fortunately there exist linear approximations of the Chezy law to ensure that the amount of
energy dissipated by bottom friction over a tidal cycles obtains the same rate as the quadratic
law. This problem was originally investigated by the Dutch physicist Lorentz. A realistic linear
approximation of the quadratic bottom drag is for instance:

Gx = ru/h (16.34)
Gy = rv/h (16.35)

where r is a properly chosen constant (typically r=0.0013). Lorentz assumed that the linear
and quadratic drag laws have to match, ie. predict the same loss of energy over 1 tidal cycle.
Lorentz worked out this problem for the M2 tide in the Waddenzee.

16.5 Linear and non-linear tides

We will summarize the consequences of non-linear acceleration terms that appear in the Laplace
tidal equations:

Linear ocean tides follow from the solution of the Laplace tidal equations whereby all forc-
ing terms, dissipative terms and friction terms can be approximated as harmonic functions.
The solution has to fulfill the condition posed by the Helmholtz equation, meaning that
the tides become a wave solution that satisfies the boundary conditions of the Helmholtz
equation. Essentially this means that ocean tides forced at a frequency resultin a
membrane solution oscillating at frequency . The surface speed of the tide is then gH.

200
 Non-linear ocean tides occur when there are significant deviations from a linear approx-
imation of the bottom drag law, or when the tide is forced through its basin geometry
along the shore or through a channel. In this case advection and bottom friction are the
main causes for the generation of so-called parasitic frequencies which manifest themselves
as undertones, overtones or cross-products of the linear tide. Examples of non-linear tides
are for instance M0 and M4 which are the result of an advective term acting on M2 . Some
examples of cross-products are M S0 and M S4 which are compound tides as a result of M2
and S2 .

16.6 Dispersion relation

Another way to look at the tide problem (or in fact many other wave problems in physics) is
to study a dispersion relation. We will do this for the simplest case in order to demonstrate
another basic property of ocean tides, namely that the decrease in the surface speed c causes
a shortening of length scale of the wave. For the dispersion relation we assume an unforced or
free wave of the following form:

u(x, y, t) = ue( j(t kx ly)) (16.36)

(
v(x, y, t) = ve j(t kx ly)) (16.37)
(
(x, y, t) = e j(t kx ly)) (16.38)

which is only defined for a local region. This generic solution is that of a surface wave, is
the angular velocity of the tide, and k and l are wave numbers that provide length scale and
direction of the wave.
To derive the dispersion relation we ignore the right hand side of eq. (16.31) and substitute
characteristic wave functions. This substitution results in:

( 2 f 2 ) = c2 k 2 + l2


(16.39)

which is a surprisingly simple relation showing that k 2 + l2 has to increase when c decreases
and visa versa. In other words, now we have shown that tidal wave lengths become shorter in
shallow waters. The effect is demonstrated in figure 16.6 with a map of the tidal amplitudes
and phases of the M2 tide in the North Sea basin.
But, there are more hidden features in the dispersion relation. The right hand side of
equation (16.39) is always positive since we only see squares of c, k and l. The left hand side is
only valid when is greater than f . Please remember that the Coriolis parameter f = 2 sin
is latitude dependent with zero at the equator. Near the equator we will always get free waves
passing from west to east or visa versa.
For frequencies equal to f one expects that there is a latitude band inside which the free
wave may exist. A nice example is the K1 tidal wave which is a dominant diurnal tide with a
period of 23 hours and 56 minutes, so that = . The conclusion is that free waves at the K1
frequency can only exist when sin is less than 1/2 which is true for a latitudes between 30N
and 30S.

201
Figure 16.6: North Sea M2 tide

202
16.7 Exercises
What is the magnitude of the Coriolis effect for a ship sailing southward at 50N with a
speed of 20 knots

Is water flowing from your tap into the kitchen sink turbulent?

What is the magnitude of a height gradient of a river with a flow of 0.5 m/s and a Chezy
coefficient of 30. The mean depth of the river is 5 meter.

What latitude extremes can we expect for free tidal waves at the Mm frequency?

How much later is the tide at Firth of Worth compared to The Wash?

What extra terms appear in the Helmholtz equation for a linear bottom drag model.

Show that advection can be written as uu

Shows that vorticity is conserved in fluid mechanics problems that are free of friction.

203
Chapter 17

Data analysis methods

Deep ocean tides are known to respond at frequencies identical to the Doodson numbers in
tables 14.3 and 14.4. Non-linearities and friction in general do cause overtones and mixed tides,
but, this effect will only appear in shallow waters or at the boundary of the domain. In the
deep oceans it is very unlikely that such effects dominate in the dynamical equations. Starting
with the property of the tides we present two well known data analysis methods used in tidal
research.

17.1 Harmonic Analysis methods

A perhaps unexpected consequence of the tidal harmonics table is that at least 18.61 years of
data would be required to separate two neighboring frequencies because of the fact that main
lines in the spectrum are modulated by smaller, but significant, side-lines. Compare for instance
table 14.3 and 14.4 where one can see that most spectral lines require at least 18.61 years
of observation data in order to separate them from side-lines. Fortunately, extensive analysis
conducted by [8] have shown that a smooth response of the sea level is likely. Therefore the
more practical approach is to take at least two Doodson numbers and to form an expression
where only a year worth of observations determine amplitude and phase of a constituent.
However, this is only possible if one assumes a fixed amplitude ratio of a side-line with respect
to a main-line where the ratio itself can be taken from the table of tidal harmonics.
Consider for instance table 14.4 where M2 is dominated by spectral lines at the Dood-
son numbers 255.555 and 255.545 and where the ratio of the amplitudes is approximately
0.02358/0.63194 = 0.03731. We will now seek an expression to model the M2 constituent:

M2 (t) = CM2 [cos(21 t M2 ) + cos(21 t + 5 t M2 )] (17.1)

where CM2 and M2 represent the amplitude and phase of the M2 tide and where = 0.03731.
Starting with:

M2 (t) = CM2 cos(21 t M2 )

+ CM2 {cos(21 t M2 ) cos(5 t) sin(21 t M2 ) sin(5 t)}

we arrive at:

M2 (t) = CM2 {(1 + cos(5 t)) cos(21 t M2 ) sin(5 t) sin(21 t M2 )} (17.2)

204
which we will write as:

M2 (t) = CM2 f (t) {cos(u(t)) cos(21 t M2 ) sin(u(t)) sin(21 t M2 )} (17.3)

or
M2 (t) = CM2 f (t) cos(21 t + u(t) M2 ) (17.4)
so that:
M2 (t) = AM2 f (t) cos(21 t + u(t)) + BM2 f (t) sin(21 t + u(t)) (17.5)
where

AM 2 = CM2 cos(M2 )
B M2 = CM2 sin(M2 )

In literature the terms AM2 and BM2 are called in-phase and quadrature or out-of-phase
coefficients of a tidal constituent, whereas the f (t) and u(t) coefficients are known as nodal
modulation factors, stemming from the fact that 5 t corresponds to the right ascension of the
ascending node of the lunar orbit. In order to get convenient equations we work out the following
system of equations: ( = 5 t):
1/2
(1 + cos())2 + ( sin())2

f (t) =
 
sin()
u(t) = arctan
1 + cos()

Finally a Taylor series around = 0 gives:

1 1 1 1
f (t) = (1 + 2 + 4 ) + ( 3 5 ) cos
4 64 8 64
1 2 1 4 1 3 5 5
+ ( + ) cos(2) + ( ) cos(3) (17.6)
4 16 8 128
75
 
5
cos(4) + cos(5) + O(6 )
64 128
1 1
u(t) = sin() 2 sin(2) + 3 sin(3)
2 3
1 4 1 5
sin(4) + sin(5) + O(6 ) (17.7)
4 5
Since is small it is possible to truncate these series at the quadratic term. The equations show
that f (t) and u(t) are only slowly varying and that they only need to be computed once when
e.g. working with a year worth of tide gauge data.
The Taylor series for the above mentioned nodal modulation factors were derived by means
of the Maple software package and approximate the more exact expressions for f and u. However
the technique seems to fail whenever increased ratios of the main line to the side line occur as
is the case with the e.g. the K2 constituent or whenever there are more side lines. A better
way of finding the nodal modulation factors is then to numerically compute at sufficiently dense
steps the values of the tide generating potential for a particular constituent at an arbitrary
locationP on Earth over the full nodalPcycle and to numerically estimate Fourier expressions like
f () = n fn cos(n.) and u() = n un sin(n.) with eq. (17.4) as a point of reference.

205
17.2 Response method
The findings of [8] indicate that ocean tides (t) can be predicted as a convolution of a smooth
weight function and the tide generating potential U a :
X
(t) = w(s)U a (t s ) (17.8)
s

with the weights w determined so that the prediction error (t) (t) is a minimum in the least
squares sense. The weights w(s) have a simple physical interpretation: they represent the sea
level response at the port (read: point of observation) to a unit impulse U a (t) = (t), hence the
name response method. The actual input function U a (t) may be regarded as a sequence of
such impulses. The scheme used in [8] is to expand U a (t) in spherical harmonics,
N X
X n
U a (, ; t) = g [anm (t)Unm (, ) + bnm (t)Vnm (, )] (17.9)
n=0 m=0

containing the complex spherical harmonics:

 1/2  1/2
m 2n + 1 (n m)!
Unm + jVnm = Ynm = (1) Pnm (cos )e( jm) (17.10)
4 (n + m)!

and to compute the coefficients anm (t) and bnm (t) for the desired time interval. The convergence
of the spherical harmonics is rapid and just a few terms n, m will do. The m-values separate
input functions according to species and the prediction formalism is:
XX
(t) = [unm (s)anm (t s ) + vnm (s)bnm (t s )] (17.11)
n,m s

where the prediction weights wnm (s) = unm (s) + jvnm (s) are determined by least-squares meth-
ods, and tabulated for each port (these take the place of the tabulated Ck and k in the harmonic
method). For each year the global tide function cnm (t) = anm (t) + jbnm (t) is computed and
the tides then predicted by forming weighted sums of c using the weights w appropriate to each
port. The spectra of the numerically generated time series c(t) have all the complexity of the
Darwin-Doodson expansion; but there is no need for carrying out this expansion, as the series
c(t) serves as direct input into the convolution prediction. There is no need to set a lower bound
on spectral lines; all lines are taken into account in an optimum sense. There is no need for
the f, u factors, for the nodal variations (and even the 20926 y variation) is already built into
c(t). In this way the response method makes explicit and general what the harmonic method
does anyway in the process of applying the f, u factors. The response method leads to a more
systematic procedure, better adapted to computer use. According to [8] its formalism is readily
extended to include nonlinear, and perhaps even meteorological effects.

17.3 Exercises
1. Why is the response method for tidal analysis more useful and successful than the harmonic
tidal analysis method, ie. what do we learn from this method what couldnt be seen with
the harmonic tide analysis method.

206
2. Design a flow diagram for a program that solves tidal amplitudes and phases from a dataset
of tide gauge readings that contains gaps and biases. Basic linear algebra operations such
as a matrix inversion should not be worked out in this flow diagram.

3. How could you see from historic tide constants at a gauge that the local morphology has
changed over time near the tide gauge.

207
Chapter 18

Load tides

Any tide in the ocean will load the sea floor which is not a rigid body. One additional meter of
water will cause 1000 kg of mass per square meter; integrated over a 100 by 100 km sea we are
suddenly dealing 1013 kg which is a lot of mass resting on the sea floor. Loading is a geophysical
phenomenon that is not unique to tides, any mass that rests on the lithosphere will cause a
loading effect. Atmospheric pressure variations, rainfall, melting of land ice and evaporation of
lakes cause similar phenomena. An important difference is whether we are dealing with a visco-
eleastic or just an elastic process. This discussion is mostly related to the time scales at which
the phenomenon is considered. For tides we only deal with elastic loading. The consequence
is that the Earths surface will deform, and that the deformation pattern extends beyond the
point where the original load occurred. In order to explain the load of a unit point mass we
introduce the Green function concept, to model the loading effect of a surface mass layer we need
a convolution model, a more efficient algorithm uses spherical harmonics, a proof is presented
in the last section of this chapter.

18.1 Green functions

In [21] it is explained that a unit mass will cause a geometric displacement at a distance from
the source:

re X 0
G() = hn Pn (cos ) (18.1)
Me
n=0

where Me is the mass of the Earth and re its radius. The Green function coefficients h0n come
from a geophysical Earth model, two versions are shown in table 18.1. The geophysical theory
from which these coefficients originate is not discussed in these lectures, instead we mention that
they represent the elastic loading effect and not the visco-elastic effect.

18.2 Loading of a surface mass layer

Ocean load tides cause vertical displacements of geodetic stations away from the load as has
been demonstrated by analysis of GPS and VLBI observations near the coast where vertical
twice daily movements can be as large as several centimeters, see for example figure 18.1. In
order to compute these maps it is necessary to compute a convolution integral where a surface

208
Figure 18.1: The top panel shows the amplitude map in millimeters of the M2 load tide, the
bottom panel shows the corresponding phase map. Note that the load tide extends beyond the
oceanic regions and that the lithosphere also deforms near the coast.

209
 Farrell Pagiatakis
n n h0n kn0 h0n kn0
1 0.1876 0.290 0 0.295 0
2 0.1126 1.001 0.308 1.007 0.309
3 0.0804 1.052 0.195 1.065 0.199
4 0.0625 1.053 0.132 1.069 0.136
5 0.0512 1.088 0.103 1.103 0.103
6 0.0433 1.147 0.089 1.164 0.093
8 0.0331 1.291 0.076 1.313 0.079
10 0.0268 1.433 0.068 1.460 0.074
18 0.0152 1.893 0.053 1.952 0.057
30 0.0092 2.320* 0.040* 2.411 0.043
50 0.0056 2.700* 0.028* 2.777 0.030
100 0.0028 3.058 0.015 3.127 0.016

Table 18.1: Factors n in equation (18.3), and the loading Love numbers computed by [21] and
by [42]. An asterisk () means that data was interpolated at n = 32, 56

mass layer, here in the form of an ocean tide chart, is multiplied times Greens functions of
angular distance from each incremental tidal load, effective up to 180 . The loading effect is
thus computed as: Z
l (, , t) = G()d M (0 , 0 , t) (18.2)

where d M represents the mass at a distance from the load. This distance is the spherical
distance between (, ) and (0 , 0 ). There is no convolution other than in and , the model
describes an instantaneous elastic response.

18.3 Computing the load tide with spherical harmonic functions

But given global definition of the ocean tide it is more convenient to express it in terms of
a sequence of load-Love numbers kn0 and h0n times the spherical harmonics of degree n of the
ocean tide. If n (, ; t) denote any nth degree spherical harmonics of the tidal height , the
secondary potential and the bottom displacement due to elastic loading are g(1 + kn0 )n n and
h0n n n respectively where:
3 w 0.563
n = = (18.3)
(2n + 1) e (2n + 1)

where w is the mean density of water and e the mean density of Earth. (Chapter 3 provides
all required mathematical background to derive the above expression, this result follows from
the convolution integral on the sphere that is evaluated with the help of spherical harmonics)
The essential difference from the formulation of the body tide is that the spherical harmonic
expansion of the ocean tide itself requires terms up to very high degree n, for adequate definition.
Farrells (1972) calculations of the load Love numbers, based on the Gutenberg-Bullen Earth

210
model, are frequently used. Table 18.1 is taken from [6] and lists a selection of both Farrells
numbers and those from a more advanced calculation by [42], based on the PREM model.
Why is it so efficient to consider a spherical harmonic development of the ocean tide maps?
Here we refer to the in-phase or quadrature components of the tide which are both treated in
the same way. The reason is that convolution integrals in the spatial domain can be solved by
multiplication of Green functions coefficients and spherical harmonic coefficients in the spectral
domain. The in-phase or quadrature ocean load tide maps contained in H(, ) follow then from
a convolution on the sphere of the Green function G() and an in-phase or quadrature ocean
tide height function contained in F (, ), for details see chapter 3.

18.4 Exercises
1. Explain how you would compute the self attraction tide signal provided that the ocean
tide signal is provided.

2. How do you compute the vertical geometric load at the center of a cylinder with a radius
of degrees.

3. Design a Green function to correct observed gravity values for the presence of mountains
and valleys, i.e. that corrects for a terrain effect. Implement this Green function in a
method that applies the correction.

211
Chapter 19

Altimetry and tides

19.1 Introduction
Satellite altimetry is nowadays an accurate technique whereby height profiles are measured
along satellite tracks over the ocean. Repeated measurement of these height profiles followed
by a suitable data analysis method provides in principle estimates of the altimetric tide. One
problem is that an altimeter will observe the sum of the solid Earth tide, an oceanic tide and
a load tide. The solid Earth tide can be modelled when the Love numbers hn are provided.
Separating the load tide from the ocean tide requires one to solve an integral equation. In this
chapter we will discuss both issues.

19.2 Aliasing
Tides observed by a satellite altimeter are usually systematically under sampled. The under
sampled diurnal and semi-diurnal frequencies result in alias periods significantly longer than the
natural periods of the tides. Any altimeter satellite has been plagued by this problem, SEASATs
lifetime (NASA altimeter, 1978) was too short for doing any serious tidal analysis, GEOSAT
(US Navy altimeter, 1985-1990) had several problems among which that the M2 tide aliases to
a period of about a year and finally ERS-1 (ESA altimeter 1991-1996) is by definition not suited
for tidal research because the sun-synchronous orbit causes all solar tides to be sampled at the
same phase.

19.3 Separating ocean tide and load tides

A satellite altimeter will observe the sum of an ocean and a load tide, where the latter is obtained
by convolution with respect to the ocean tide, thus we have:

Sa = So + L(So ) (19.1)

where Sa is the tide observed by the altimeter, and where So is a ocean tide. Operator L() is a
convolution integral as explained in chapter 18. In order to obtain ocean and load tides we have
to solve an integral equation. Since L is a linear operator the ocean tide is obtained by:

So = (I + L)1 Sa (19.2)

212
 Authors version Q1 O1 P1 K1 N2 M2 S2 K2
Schwiderski 1980 0.34 1.23 0.61 1.44 1.19 3.84 1.66 0.59
Cartwright-Ray 1991 1.22 0.63 1.89 0.96 3.23 2.22
Le Provost et al. meom94.1 0.28 1.04 0.46 1.23 0.87 2.99 1.56 0.50
Egbert et al. tpxo.1 0.96 1.26 2.30 1.55
Egbert et al. tpxo.2 0.29 0.98 0.45 1.32 0.76 2.27 1.26 0.56
Sanchez-Pavlis gsfc94a 0.35 1.06 0.54 1.41 0.86 2.31 1.23 0.66
Ray et al. 1994 0.37 1.00 0.40 1.25 0.81 2.04 1.23 0.51
Schrama-Ray 1993.10 1.15 1.35 2.02 1.26
Schrama-Ray 1994.11 1.02 1.19 0.85 1.85 1.20

Table 19.1: Ground truth comparison at 102 tide gauges, the first two tide models are developed
before T/P. Le Provost et al. ran a global finite element model that is free from T/P data.
Egbert et al., also ran a finite element model while assimilating T/P data. Sanchez & Pavlis
and Ray et al. used so-called Proudman functions to model the tides, they did incorporate T/P
data. Schrama & Ray applied a straightforward harmonic analysis to the T/P data to determine
improvements with respect to a number of tidal constituents.

It turns out that there is a fast inversion algorithm capable of inverting this problem within
several iterations
(0)
Sl = L(Sa )
(0)
So(0) = Sa Sl
(1)
Sl = L(So(0) )
(1)
So(1) = Sa Sl
(2)
Sl = L(So(1) )
(2)
So(2) = Sa Sl
..
.

This procedure has been used to separate the ocean and load tide from TOPEX/POSEIDON
(T/P) altimetry data.

19.4 Results
To close this chapter on tides we want to mention that the T/P satellite altimeter mission
(NASA/CNES, active since August 1993) has stimulated the development of a series of new tide
models more accurate than any previous global hydrodynamic model, see for instance [56]. The
main reason for the success of the T/P mission in modeling the deep ocean tides should be seen
in the context of the design of the mission where the choice of the nominal orbit is such that
all main tidal constituents alias to relatively short periods. A few of the results are tabulated
in table 19.1 where the r.m.s. comparisons to 102 ground-truth stations in (cm) are shown.
Ocean tides in shallow coastal areas are not that easily observed with T/P altimetry because
of the non-harmonic response of tides in shallow seas leading to spatial details exceeding the

213
resolution attainable by the T/P inter track spacing. This behavior was explained in chapter 16,
in particular at the point where the dispersion relation of barotropic waves was discussed. For
shallow seas it is in general better to rely on regional tide/storm surge models. An example
for the North Sea area is the Continental Shelf Model (CSM) maintained by the RIKZ group,
Department of Rijkswaterstaat, Koningskade 4, 2500 EX Den Haag, The Netherlands.

19.5 Exercises
1. Show that the recursive algorithm to solve eq. (19.2) is valid.

2. What is the aliasing period of the M2 tide when it is observed from the Envisat orbit
which is a 35 day sun-synchronous repeat orbit. Can you also observe the S2 tide with an
altimeter from this orbit?

3. The T/P orbit completes 127 orbital periods in 10 nodal days. Use the J2 gravity preces-
sion equations to find the proper orbital altitude at an inclination of 66 degrees and an
eccentricity of 0.001. What is the ground track repeat time.

4. Use the answers of the previous question to compute the aliasing period of the M2 and
the S2 tide.

5. How much time does it take to disentangle Ssa and K1 from T/P.

214
Chapter 20

Tidal Energy Dissipation

20.1 Introduction
This chapter is about tidal energy computations in the Earth-Moon system. The subject is
known for quite some time, a comprehensive reference can be found in [31] where tidal energetics
is described prior to the refinement of tidal models by satellite altimetry, and in particular from
the T/P mission, see [56] and [49]. Tidal energy dissipation discusses the way how mechanical
energy stored in tidal motions is converted into another form of energy. Where this process
actually occurs and into which form energy is converted are separate questions that we will
discuss later in this chapter. Basic observations confirming that energy is dissipated in oceanic
tides are linked to the slowdown of Earth rotation, which is about 5 1022 rad/s2 , and
lengthening of the distance between the Earth and the Moon by about 3.82 0.07 cm/year,
see also [17], [40], [39] and [31]. To explain this mechanism we will review the Earth-Moon
configuration like shown in figure 20.1: According to [17] the global rate of energy dissipation
is 2.50 0.05 Terawatts (TW) which is relevant for the M2 tide. In section 20.2 we will look
closer into the method employed by [17] which is entirely based on the LLR observation that
the semi-major axis of the lunar orbit increases over time. From this information alone one can
reconstruct the rate of dissipation of the M2 tide globally. Since 1969 satellite altimetry has
opened alternative ways to obtain estimates of the rate of tidal energy dissipation. The reason is
that the shape of the tidal ellipsoid that is pointing to the Moon (and the Sun) can be measured
directly. This allows us to compare the LLR method to the satellite altimetry, and to identify
where dissipation occurs within the Earth system on a global scale. Once the tides are mapped
in detail in the ocean, we can go even one step further, the dissipation estimates can be refined
to a local level.
The problem of sketching a complete picture of the dissipation mechanisms is clearly a
multidisciplinary scientific challenge where astronomy, geodesy, physical oceanography and me-
teorology come together. Purpose of writing this chapter is to go through the derivation of the
tidal energy equations and to confirm the global dissipation rates in the oceanic tides from a
handful of existing satellite altimetry ocean tide models. For this purpose dissipation in ocean
tide models is treated from a fluid dynamic point of view which is discussed in section 20.3.
The evaluation of global tidal energy dissipation problem based upon tide models obtained from
satellite altimetry is discussed in section 20.4.

215
Figure 20.1: This figure shows how the Earth spin rate slows down as a result of the gravitational
torque formed by the acceleration vectors a and b. The Moon is also slowed down in this
configuration, causing it to move away from Earth at a rate of 3.82 0.07 cm per year which is
observed by lunar laser ranging (LLR) to reflector packages left on the lunar surface since 1969.

Figure 20.2: Left: Lunar laser ranging instrument in action at the Goddard Space Flight Center,
Greenbelt Maryland, USA. Right: laser reflector at the lunar surface installed during the Apollo
11 mission in 1969. Images: NASA

216
20.2 Tidal energetics from lunar laser ranging
20.2.1 Introduction
Tidal energy dissipation is the process whereby mechanical energy is converted into any other
form of energy such as heat as a result of friction, or the energy required to mix fluids of different
densities. In subsection 20.2.2 the set-up of the problem is discussed, in subsection 20.2.3 the
global energy dissipation for M2 is derived from the recession of the lunar semi-major axis
observed by lunar laser ranging (LLR). In subsection 20.2.4 we explain where M2 dissipation
occurs in the Earth-Moon system and how it relates to the 2.42 TW for M2 observed by satellite
altimetry.

20.2.2 Relevant observations

Presently we see one face of the Moon in its orbit about the Earth and we speak about a
phase-locked configuration. The Moon itself has dissipated all its rotational energy, and a tidal
ellipsoid on the Moon today would be oriented with its main bulge pointing towards Earth.The
Moon is assumed to follow a Kepler orbit about the Earth, to compute this orbit you need the
gravitational constant of the Earth e and that of the Moon m . The semi-major axis of the
lunar orbit is known, it follows in fact from observations, we call it a whereby a = 3.84399 108
m. We also know from lunar orbit analysis the rate of change of a which we call a. The latter is
obtained by laser ranging to reflector packages left on the lunar surface in the 70s, see figure 20.2,
and this resulted in the already advertised value of a of 3.82 0.07 cm per year. Other constants
(estimated from observations) are e = 3.9860044 1014 m3 s2 , and m = 4.9048695 1012
m3 s2 , an Earth year lasts 365.24 days, and the sidereal day length is 86400 (365.241)
365.24 86164
seconds. The sidereal day is the time in seconds we need to see the same star crossing a local
meridian twice, presently we use other observation techniques such as VLBI, GPS and satellite
laser ranging to determine the length of a sidereal day. Also we need the dimensions of the
Earth and Moon, just a spherical approximation is sufficient, re = 6378137 m and rm = 1737.1
km are the radii of Earth and Moon respectively. Finally there is the universal gravitational
constant G which is measured in the laboratory, it is measured for instance with a torsion
balance and its value is G = 6.6740831 1011 Nm2 kg2 . These are for now the only constants
required to solve the M2 dissipation problem with the LLR method. Crucial for understanding
the astronomical dissipation at M2 is of course the rate of change of the lunar orbit a, this turns
out to be the only value you need to solve the dissipation problem for M2 , the relative error
of a is significantly greater than all other constants involved in the calculation, it sets an error
bound on the estimated dissipation for M2 .

20.2.3 Modeling
Tidal energy dissipation in the Earth-Moon system results directly in an increase of the semi
major axis a of the lunar orbit. Responsible for a is the oceanic tidal bulge which is misaligned
with the Earth-Moon vector. The tidal bulge is modeled with an amplitude and a phase lag,
which are both observed by satellite altimetry as will be discussed in section 20.3. The phase
lag of the tidal bulge on Earth is positive because fluids need extra time to respond to the
tidal forcing which sets the ocean (and solid earth and atmospheric tides) in motion. The tidal
amplitude and phase lag provide sufficient information on how to compute the tidal torque that

217
will slow down Earth rotation. Crucial for the problem is that Earth rotation has the same
sense of rotation as the Moon in orbit about the Earth. Also crucial for the problem is that the
eigen-rotation of the Moon is such that it is presently phase-locked with Earth rotation.
There are two ways to compute tidal energy dissipation. Either we can make use of the
observed amplitude and phase lags of the tidal bulge, or we make use of the lunar laser ranging
(LLR) method where a is directly observed. We will proceed in the following with the LLR
method which is sometimes also referred to as the astronomic method. The energy equation
that describes the configuration is as follows:
E
D= (20.1)
t
where D is the dissipation in the system, the units are in Watt (read Joules per second) and E
is the energy (or work, read Joules) in the problem. The work term breaks down in three parts:

(e + m )mm 1 1
E = E1 + E2 + E3 = + Je e2 + Jm m
2
(20.2)
2a 2 2
Here E1 is the energy required to keep the Moon in orbit about the Earth, this equation came
from equation (2.22) where we added a kinetic and potential, next we scaled these problem to
account for the lunar mass mm . The E2 term describes the energy stored in Earth rotation,
likewise E3 is the rotational energy stored in the Moons eigen-rotation. For the moments of
inertia you can assume a rotating sphere, in other words, Je = 25 me re2 and Jm = 25 mm rm
2 . We

also need the mass of the Earth me and the mass of the Moon, you get them by dividing the
gravitational constants (which are observed) by the universal gravitational constant. Let me
now piece by piece discuss how to proceed to get D. The first term is:
E1 e + m
D1 = = mm a (20.3)
t 2a2
which describes the power needed to move the Moon to a higher orbit. The second term
D2 = Et requires knowledge of the eigen-rotation of the moon and the time derivative of eigen-
2

rotation. Both terms follow (due to the phase-locked configuration) from the mean motion of
the lunar orbit:
e + m 1/2
 
3 G(mm + me )
m = 3
m = a (20.4)
a 2 m a4
To obtain the slow-down of Earth rotation we make use of the fact that the tidal bulge on Earth
is phase-locked with the lunar orbit. The rate of Earth rotation e and the mean motion of
the lunar orbit m result together in Ld which is the lunar period in siderial days. With Ld we
compute e and the dissipation term D2 :
e
Ld = e = m Ld D2 = Je e e (20.5)
m
For the last term in the dissipation equation we obtain, this term describes the slowdown of the
lunar eigen-rotation:
E3
D3 = = Jm m m (20.6)
t
The numerical values for all terms are as follows: D1 = 0.121 TW, D2 = 2.441 TW, and
D3 = 2.977 MW, (TW stands for TeraWatt or 1012 Watt, MW stands for MegaWatt or 106

218
Watt). The sum of all terms is D = 2.563 0.047 TW which is called the astronomic value
for the rate of energy dissipation of the M2 tide, its confidence interval follows mostly from the
observed values for a.
With the help of a maple program you can assemble all contributing terms and obtain a
direct expression for the dissipation in the Earth-Moon system:

D = f (a, m, M, rm , re , e , G) a
 
1 2 6 2 2 2 3
f () = 5mm (me + mm )(a + rm )G + 6me re e a
10a4 5

from which we conclude that the f () term only depends on the current configuration of the Earth-
Moon mass ratios, their radii, Earth rotation rate, and the lunar semi-major axis. Dissipation
at M2 is therefore measurable by observing the rate of the recession a of the lunar semi-major
axis.

Earth-Moon system in the past

It is tempting to use the results of the LLR estimate for tidal energy dissipation to reconstruct
the Earth-Moon system before the present day, see also [31]. We found expressions to reconstruct
a and e as a function of D and we could integrate backward in time. The bottleneck in this
discussion is the behavior of D in the past, because this is a term that depends on the average
depth of the oceans and the abundance of continental shelves where most of the dissipation
takes place. If the oceans in the past had many continental shelves, or, if the ocean basins were
shaped such that resonance took place, then D would certainly be different compared to present
day situation. The results for D depend on whether the Earth-Moon system is in a phase-locked
(tidally-locked) configuration, if the system is not phase-locked then predicting the past becomes
even more difficult. Geologic survey of microbial lifeforms in the tidal pools could be used to
constrain the paleo ocean tide models, see [31] for more details.

20.2.4 Interpretation
To summarize the result of the LLR method, for all dissipation terms we find:

The largest term is D2 and this refers to the slowdown of Earth rotation, the rate of
slowdown of Earth rotation follows directly from the fact that the tidal bulge is phase
locked with the lunar orbit, also, the tidal bulge for M2 leads the lunar sub-point.

The second largest term is D1 which describes the dissipation related to increasing the
semi-major axis of the lunar orbit.

Finally there is D3 which is a minor term describing the loss of eigen-rotation of the Moon,
it may also be ignored for this problem.

The only significant uncertainty in this calculation of D is the confidence interval of a, this
is approximately 2% of the observed rate of recession of the Moon, it is the only significant
uncertainty for the rate of energy dissipation at M2 observed by LLR.

The global rate of energy dissipation D for the M2 tide obtained from LLR can be compared
to independent values obtained from satellite altimetry where we find 2.42 TW for M2 . (A

219
discussion of the altimeter results follows later in section 20.3.) As a result there is a difference
of 0.12 TW between altimetry and LLR which is too large because of the uncertainty limit of
the LLR method. More important is that there is a physical cause to explain this difference,
satellite altimetry will be sensitive to dissipation in the ocean, and it will not see a solid Earth
dissipation, but the LLR method does pick this up. For this reason 0.12 TW is thought to be
dissipating in the solid Earth cf. [50]. An independent observation of solid Earth dissipation at
M2 does not really exist, terrestrial gravimetry would be a suitable technique (theoretically) but
the accuracy of terrestrial gravimetry is not sufficient to confirm an amplitude and phase lag of
the solid-earth body tide.

20.3 Tidal energetics and fluid dynamics

We start with the equations of a fluid in motion and show the necessary steps to arrive at the
energy equation which contains a work term, a divergence term and a dissipation term. We will
integrate this equation over a tidal cycle and over the oceans to confirm that the dissipation
term equals the work term. In an example we demonstrate that the global dissipation rate at
M2 is 2.41 TW for the GOT99.2 model cf. [49]. The dissipation rates at other constituents such
as O1 K1 and S2 are smaller; they are respectively 0.17, 0.33 and 0.43 TW.

20.3.1 Dissipation terms in the Laplace Tidal equations

We start with the equations of motion whereby the velocity terms u are averaged over a water
column, see also [23] or within these notes eq. (16.23)(a-c):

t u + f u = g + F (20.7)
t = . (uH) (20.8)

In these equations H is the height of the water column, is the surface elevation, f is the
Coriolis vector, g is the gravitational acceleration, is the acceleration term that sets water in
motion and F contains terms that model the dissipation of energy or terms that model advection.
Essentially the momentum equations (20.7) state that the Coriolis effect, local gravity and the
gradient of the pressure field are balanced while the continuity equation (20.8) enforces that
there are no additional drains and sources.
For tidal problems the forcing function is a summation of harmonic functions depending on
indicating the frequency of a tidal line. If F is linear, in the sense that we dont allow squaring
of u and , while imposing harmonic boundary conditions at frequency then solutions for u
and will also take place at . However if F contains advective or non-linear frictional terms
both causing a velocity squaring effect then the equations become non-linear so that solutions of
u and will contain other frequencies being the sums of differences of individual tidal lines. By
means of scaling considerations one can show, see [6], that non-linearities only play a marginal
role and that they are only significant in coastal seas. An example is the overtone of M2 (called
M4 ) which is small in the open oceans, see also chapter 16.
In [4] we find that the energy equation is obtained by multiplying the momentum equations
(20.7) times Hu and the continuity equation (20.8) times g with representing the mean
density of sea water. (Unless it is mentioned otherwise we assume that = w ). As a result we

220
obtain:  
1 1
t H(u2 + v 2 ) + g 2 = gH. (u) + Hu. Hu.F (20.9)
2 2
where we used the property (ab) = ab + ba. In the following we evaluate the time average
over a tidal period by integrating all terms in eq. (20.9) over a tidal period T where T = 2/.
In order to condense typesetting a new notation is introduced:
1 t=T +c
Z
<F >= F (t) dt
T t=c
where we remind that:  
1 1
< t H(u2 + v 2 ) + g 2 >= 0
2 2
due to the fact that u = (u, v) and are harmonic functions. (Note: formally the continuity
equation should contain a term H + instead of just H, yet  H so that the effect can be
ignored in the computations.) Characteristic in the discussion of the energy equation is that
averaging will not cancel the remaining terms in eq. (20.9). We obtain:
<W >+<P >=<D> (20.10)
where < W > is the gravitational input or work put into the tides:
< W > = H < u. >
with < P > denoting the divergence of energy flux with:
< P > = gH. < u >
The dissipation of energy < D > is entirely due to F :
< D > = H < u.F >
To obtain the rate at which tidal energy is dissipated eq. (20.10) should be integrated locally
over a patch of ocean or globally over the entire oceanic domain, see also [4] [6] [16] [31] [39].
The results will be discussed later in these lecture notes.

20.3.2 A different formulation of the energy equation

Let be the oceanic tide, e the equilibrium tide and sal the self-attraction and loading tide
and U the volume transport then, cf. [16]:
< D >= g. < U > +g < U e > +g < U sal >
where U = Hu and X
e = g 1 (1 + kn hn )Una
n
with Una denoting the astronomical tide potential and hn and kn Love numbers for the geometric
radial deformation and the induced potential that accompanies this deformation. The self-
attraction and loading tide sal is:
X 3(w /e )
sal = g 1 (1 + kn0 h0n ) nma Ynma (, )
nma
(2n + 1)

221
where e is the mean density of the Earth while h0n and kn0 are load Love numbers. In this equa-
tion nma are spherical harmonic coefficients of the ocean tide and Ynma (, ) denote spherical
harmonic functions. To avoid confusion we mention that our normalization terms are chosen
such that: Z
2
Ynma (, ) d = 4

where 
cos(m)P nm (cos ) : a = 0
Ynma (, ) =
sin(m)P nm (cos ) : a = 1
where and denote geographic longitude and co-latitude.

20.3.3 Integration over a surface

So far equation (20.10) applies to a local patch of ocean. If we are interested in a dissipation
rate over a domain then it is necessary to evaluate the surface integral. For the work integral
we can use the property:
Z Z Z
<W > =
\ H < u. > d = < H.(u) > d < H.u > d (20.11)

where the continuity equation .(uH) = t is applied. After integrating all terms we get:

<\
W1 > + <\
W2 > + <\
P > = <\
D> (20.12)
where: Z

<\
W1 > = < > d (20.13)
t
Z
< W2 > =
\ < .(Hu) > d (20.14)

Z
<P > =
\ < g.(Hu) > d (20.15)

For completeness it should be mentioned that the surface integrals for <\
W2 > and <\P > may
be replaced by line integrals over an element ds along the boundary of , cf. [4]:
I
< W2 > =
\ < H(u.n) > ds (20.16)

and I
<\
P > = < g H(u.n) > ds (20.17)

where n is a vector perpendicular to .

20.3.4 Global rate of energy dissipation

In case our integration domain concerns the global domain we can assume that <\ W2 > = 0
and <\ P > = 0 since the corresponding surface integrals can be written as line integrals along
the boundary where we know that the condition (u.n) = 0 applies. The conclusion is that
the global dissipation rate can be derived by <\
D > = <\ W1 >, meaning that we only require
knowledge of the function and the ocean tide .

222
Spherical harmonics
At this point it is convenient to switch to spherical harmonic representations of all relevant
terms that are integrated in the work integral because of orthogonality properties, see also [31].
A convenient representation of the oceanic tide field is a series of global grids whereby an
in-phase and a quadrature version are provided for a selected number of constituents in the
diurnal and semi-diurnal frequency band. The problem of representing can be found in [6]
where it is shown that:
X
(, , t) = f [P (, ) cos((t) u ) + Q (, ) sin((t) u )] (20.18)

The definitions of f and u are related to the effect of side lines modulating the main wave,
see also section 17.1. In the following discussion we will ignore the effect of f and u (ie.
f = 1 and u = 0) and assume that their contribution can be neglected in the evaluation of
the energy equation. In essence this assumption says that we convert the formal definition of a
tidal constituent into that of a single wave at frequency .

Prograde and retrograde waves

To appreciate the physics of tidal energy dissipation [31] presents a wave splitting method. The
essence of this method is that we get prograde and retrograde waves which are constructed
from the spherical harmonic coefficients of P and Q in eq. (20.18) at a given frequency . To
retrieve both wave types we develop P and Q in spherical harmonics:
X
P = [anm cos m + bnm sin m] P nm (cos ) (20.19)
nm
X
Q = [cnm cos m + dnm sin m] P nm (cos ) (20.20)
nm

to arrive at:
X
+ +

(, , t) = Dnm cos((t) + m nm ) + Dnm cos((t) m nm ) P nm (cos ) (20.21)
nm

with:
1
Dnm cos(nm ) = (anm dnm ) (20.22)
2
1
Dnm sin(nm ) = (cnm bnm ) (20.23)
2
In this notation the wave selected with the + sign is prograde; it is a phase locked wave that
leads the astronomical bulge with a certain phase lag. The second solution indicated with the
sign is a retrograde wave that will be ignored in further computations. From here on Dnm +
+ are the only components that remain in the global work integral < W
and nm c1 >.
Tables of spherical harmonic coefficients and associated prograde and retrograde amplitudes
and phase lags exist for several ocean tide solutions, see also [49] who provides tables of 4 diurnal
waves Q1 O1 P1 K1 and 4 semi-diurnal waves N2 M2 S2 K2 . The required Dnm and terms
nm
are directly derived from the above equations, albeit that our spherical harmonic coefficients
bnm and dnm come with a negative sign compared to [49].

223
Analytical expression for the global rate of dissipation
In the following we will apply the coefficients anm through dnm in eqns. (20.19) and (20.20)
in the evaluation of eq.(20.13). We require the time derivative of the tidal heights and the
function, a discussion of both terms and their substitution in eq.(20.13) is shown hereafter.

Forcing function
For the forcing function we know that it is directly related to the astronomical tide generation
function Una and secondary potentials that follow from the self attraction and loading tide:

= g (e + sal ) (20.24)

However from this point on we concentrate of the e term assuming that the sal term is smaller.
The justification for assuming = g e is that an equilibrium ocean tide should be achieved in
case there are no tidal currents u and terms F , see also eq. (20.7). In addition we know from
[6] that for all dominant tidal waves we always deal with n = 2 and m = 1 for the diurnal cases
and m = 2 for the semi-diurnal cases. According to [6] the expression for U2a for a diurnal wave
at frequency with (n + m) : odd is:
a
Un=2 = A21 P 21 (cos ) sin((t) + m) (20.25)

while the expression for U2a for a semi-diurnal wave at frequency with (n + m) : even is:
a
Un=2 = A22 P 22 (cos ) cos((t) + m) (20.26)

Time derivative of the elevation field

The t term in the < W
c1 > integral is defined on basis of the choice of where we will only
use the prograde component:
X
+ +
= Dnm sin((t) + m nm )P nm (cos ) (20.27)
t nma

Phase definitions of the ocean and the astronomical tide generating potential are both controlled
by the expression (t) and the geographic longitude . Due to the fact that we average over a
full tidal cycle T it doesnt really matter in which way (t) is defined as long as it is internally
consistent between t and .

Result
We continue with the evaluation of m = 1 for diurnal waves and m = 2 for semi-diurnal waves
and get:
+
Z  
+ cos 2m
<\D> = < > d = Wnm D2m + (20.28)
t + sin 2m
+
with Wnm = 4R2 (1+k2 h2 )A2m where R is the mean Earth radius and whereby cos 2m
+
is evaluated for the diurnal tides and the sin 2m for the semi diurnal tides. We remind that
eq. (20.28) matches eq.(4.3.16) in [31]. The diurnal equivalent does however not appear in this
reference and phase corrections of /2 should be applied. In addition we notice that we did not

224
 Q1 O1 P1 K1 N2 M2 S2 K2
SW80 0.007 0.176 0.033 0.297 0.094 1.896 0.308 0.024
FES94.1 0.007 0.174 0.035 0.321 0.097 2.324 0.350 0.027
FES95.2 0.007 0.186 0.035 0.310 0.111 2.385 0.390 0.027
FES99 0.008 0.185 0.033 0.299 0.109 2.438 0.367 0.028
SR950308 0.006 0.150 0.028 0.233 0.112 2.437 0.434 0.027
SR950308c 0.007 0.180 0.034 0.288 0.114 2.473 0.435 0.027
GOT99.2 0.008 0.181 0.032 0.286 0.110 2.414 0.428 0.029
TPXO5.1 0.008 0.186 0.032 0.293 0.110 2.409 0.376 0.030
NAO99b 0.007 0.185 0.032 0.294 0.109 2.435 0.414 0.035
CSR40 0.008 0.181 0.031 0.286 0.111 2.425 0.383 0.028
Mean 0.007 0.179 0.032 0.290 0.109 2.416 0.397 0.029
Sigma 0.001 0.012 0.002 0.024 0.005 0.042 0.031 0.002

Table 20.1: Dissipation rates of 10 tide models, the model labels are explained in the text, the
average and standard deviations are computed over all models except SW80, units: Terawatt

take into account the effect of self attraction and loading tides in the evaluation of the global
dissipation rates although this effect is probably smaller than the oceanic effect. The closed
expression for the self attraction and loading effect is:
0 h0 ) +
 
+ 3(1 + k2 2 w cos 2m
<\ D > = Wnm D2m + (20.29)
5e + sin 2m

which follows the same evaluation rules as eq.(20.28).

20.4 Rate of energy dissipation obtained from ocean tide models

We compute the global dissipation rates for eight tidal constituents which are considered to be
energetic, meaning that their harmonic coefficients stand out in the tide generating potential.
The rates corresponding to eqn. (20.28) for the diurnal constituents Q1 , O1 , P1 and K1 and the
semi-diurnal constituents N2 M2 S2 and K2 are shown in table 20.1. For we have used 1026
kg/m3 , h2 = 0.606, k2 = 0.313 and R = 6378.137 km.
The models in table 20.1 are selected as follows: 1) availability of the model, 2) its ability
to provide a global coverage of the oceans, and 3) documentation to retrieve the in-phase and
quadrature coefficient maps from the data.

20.4.1 Models
The SW80 and the FES94.1 models did not rely on altimeter data and should be seen as
hydrodynamic estimates of the ocean tides, dissipation of the model was estimated and this also
constrains so that the observed data is in agreement with the prediction from fluid dynamics
differential equations. The SW80 model is described in [57], [58] and [59] and is often referred
to as the Schwiderski model providing at its time the first realistic hydrodynamic estimate of
the ocean tides obtained by solving the Laplace tidal equations. An more modern version is the

225
FES94.1 model. It is a finite element solution (FES) with the ability to follow the details of
the tides in shallow waters. Version 94.1 is documented in the JGR Oceans special volume on
the T/P altimetry system, see [5]. The FES95.2 model is a refinement of the FES94.1 model
that relies on the representer technique described by [15] to assimilate T/P altimetry data. The
FES99 model is new version of the FES95.2 model that incorporates a larger time span of the
T/P data which comes in the form of spatially filtered altimetry data at a number of crossover
locations. The FES99 model assimilates both T/P crossover data and tide gauge data.
In table 20.1 there are four empirical tide models that heavily rely on tidal constants directly
estimated from the T/P altimeter data set. The SR950308 model is an updated version of the
method documented by [56] and is based upon a local harmonic improvement of the in-phase
and quadrature components relative to a background ocean tide model. Thereby it relies on the
availability of T/P data and not so much on model dynamics. In the above table the SR950308
model is evaluated within latitude bands that follow from the orbit inclination of T/P. The
SR950308c model is an identical version that is complemented by SW80 tidal constants outside
the range of the SR950308 model. Both the SR models are based upon cycles 2 to 71 of T/P
altimetry. Another empirical model is the GOT99.2 model that is documented in [49]. It is
based on the same technique as described in [56] and can be seen as an update to the earlier
approach in the sense that 232 TOPEX cycles are used rather than the 70 cycles available at
the time the SR950308 model was developed.
The CSR4.0 model falls essentially in the same category of methods as the SR950308 and
the GOT99.2 model. In essence it is an empirical estimation technique and an update to the
CSR3.0 model documented in cf. [13]. The CSR4.0 model is based upon an implementation
of a spectral response method that involves the computation of orthotides as described in the
paper of [24]. Spectral response models enable to take the effects of minor tidal lines into account
without separately estimating individual harmonic coefficients of those lines. Without doubt this
procedure relaxes the parameter estimation effort. A drawback of the used orthotide method
is that resonance effects or energy concentrated at tidal cusps in the tides leak to neighboring
lines.
Two other models that we included in table 20.1 are TPXO5.1 and NAO99b. The TPXO5.1
model is based upon the representer approach as described in [14] whereby T/P crossover data
is assimilated in the solution. It differs from the FES95.2 and FES99 models; the method of
discretization and dynamical modelling are set-up in different ways. The NAO99b model, cf. [28],
is also based upon a data assimilation technique. In this case a nudging technique rather than
a representer technique is used.

20.4.2 Interpretation
Table 20.1 shows that most dissipation rates of the selected tide models differ by about 2%. The
average global dissipation rate of M2 is now 2.42 TW and its standard deviation is 0.04 TW.
The SW80 and the FES94.1 models are the only two exceptions that underestimate the M2
dissipation by respectively 0.5 and 0.1 TW. In [6] it is mentioned that this behavior is typical
for most hydrodynamic models that depend (for their dissipation rates) on the prescribed drag
laws in the model. All other post T/P models handle this problem in a different way, and are
based upon assimilation techniques.
Other tidal constituents that stand out in the dissipation computations are O1 K1 and S2 .
For the latter term it should be remarked that energy is not only dissipated in the ocean, but also

226
in the atmosphere. This can be confirmed by comparing the S2 dissipation to an independent
geodetic estimate from satellite geodesy.

20.5 Local estimates of tidal energy dissipation, internal waves

and mixing
In [16] one finds an estimate of 0.7 TW of energy dissipation in the deep oceans for the M2 tide
which is mainly attributed to internal wave generation at sub-surface ridges and at continental
shelf boundaries, the relevant charts are show in figure 20.3. Note that the dissipation mechanism
differs from that on continental shelves where bottom current friction is responsible for the
dissipation, in the deep oceans bottom drag is small, and energy will dissipate in another way,
namely be mixing of light surface waters which are on top of deeper more saline ocean water.
Mixing takes energy, and this explains why numerical ocean tide models that fit to the altimeter
data also require dissipation terms in regions where mixing takes place. As a result we see
relatively large local dissipations near the Hawaiian ridge system in figure 20.3 which can not be
explained by bottom boundary friction. The relevance of energy dissipation in the deep oceans
is that mixing by internal waves is partly responsible for maintenance of the observed abyssal
density stratification. The required energy to maintain this stratification requires, according to
[41] of the order of 2 TW. Internal tides are according to [16] responsible for approximately 1
TW, 0.7 TW is confirmed for the M2 tide while the remainder comes from other tidal lines. To
bring the total up to 2 TW mentioned by [39] we need an extra 1 TW from mixing by wind.

20.6 Exercises
Why does an orbital analysis of Lageos and Starlette tracking data give us a different value
for the dissipation on S2 compared to dissipation estimates from satellite altimetry?

Is there an age limit on our solar system given the current rate tidal energy dissipation?

How would you measure the rate of energy dissipation for M2 in the North sea if transport
measurements are provided at the boundary of a model for the North sea, and if tidal
constants for are provided within the numerical box?

Verify whether tidal energy dissipation on a planet circularizes the orbit of a Moon re-
sponsible for generating the tides.

On the Jovian moon Io we also have tidal energy dissipation, the moon is in an eccentric
trajectory around Jupiter, explain with a drawing how the tidal bulge on Io moves around
in an orbit around Jupiter.

What will change in the LLR method to estimate a when a moon is within a geostationary
altitude, like is the case for Phobos orbiting Mars. Explain the energy equations.

227
Figure 20.3: Top: amplitude and phase lines of the M2 ocean tides determined from satellite
altimetry. Bottom: dissipation densities derived from a deep ocean numerical ocean model
assimilating altimeter data. Images: NASA

228
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