Introduction To Diff Erential Equations Jeff Rey R. Chasnov

Introduction to Dierential
Equations
Lecture notes for MATH 150/151
Jerey R. Chasnov
The Hong Kong University of
Science and Technology
The Hong Kong University of Science and Technology
Department of Mathematics
Clear Water Bay, Kowloon
Hong Kong
Copyright c 2009 by Jerey Robert Chasnov
All rights reserved.
Unlimited permission to copy or use is hereby granted subject to inclusion of this
copyright notice.
Preface
What follows are my lecture notes for Math 150/151: Introduction to ordi-
nary dierential equations/Dierential equations and applications, taught at
the Hong Kong University of Science and Technology. Math 150, with two lec-
ture hours per week, is primarily for non-mathematics majors and is required
by several engineering and science departments; Math 151, with three lecture
hours per week, is primarily for mathematics majors and is required for applied
mathematics students.
All web surfers are welcome to download these notes and to use them freely
for teaching and learning. I welcome any comments, suggestions, or corrections
sent to me by email (jerey.chasnov@ust.hk).
iii
Contents
0 A short mathematical review 1
0.1 The trigonometric functions . . . . . . . . . . . . . . . . . . . . . 1
0.2 The exponential function and the natural logarithm . . . . . . . 1
0.3 Denition of the derivative . . . . . . . . . . . . . . . . . . . . . 2
0.4 Dierentiating a combination of functions . . . . . . . . . . . . . 2
0.4.1 The sum or dierence rule . . . . . . . . . . . . . . . . . . 2
0.4.2 The product rule . . . . . . . . . . . . . . . . . . . . . . . 2
0.4.3 The quotient rule . . . . . . . . . . . . . . . . . . . . . . . 2
0.4.4 The chain rule . . . . . . . . . . . . . . . . . . . . . . . . 3
0.5 Dierentiating elementary functions . . . . . . . . . . . . . . . . 3
0.5.1 The power rule . . . . . . . . . . . . . . . . . . . . . . . . 3
0.5.2 Trigonometric functions . . . . . . . . . . . . . . . . . . . 3
0.5.3 Exponential and natural logarithm functions . . . . . . . 3
0.6 Denition of the integral . . . . . . . . . . . . . . . . . . . . . . . 3
0.7 The fundamental theorem of calculus . . . . . . . . . . . . . . . . 4
0.8 Denite and indenite integrals . . . . . . . . . . . . . . . . . . . 4
0.9 Indenite integrals of elementary functions . . . . . . . . . . . . . 5
0.10 Substitution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 5
0.11 Integration by parts . . . . . . . . . . . . . . . . . . . . . . . . . 6
0.12 Taylor series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6
0.13 Complex numbers . . . . . . . . . . . . . . . . . . . . . . . . . . 7
1 Introduction to odes 11
1.1 Denition of an ode . . . . . . . . . . . . . . . . . . . . . . . . . 11
1.2 The simplest type of dierential equation . . . . . . . . . . . . . 11
2 First-order odes 13
2.1 Separable equations . . . . . . . . . . . . . . . . . . . . . . . . . 13
2.2 Linear equations . . . . . . . . . . . . . . . . . . . . . . . . . . . 17
2.3 Applications . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
2.3.1 Compound interest . . . . . . . . . . . . . . . . . . . . . . 19
2.3.2 Annual percentage yield (APY) . . . . . . . . . . . . . . . 20
2.3.3 The rule of 72 . . . . . . . . . . . . . . . . . . . . . . . . . 20
2.3.4 Compound interest with deposits or withdrawals . . . . . 21
2.3.5 Chemical reactions . . . . . . . . . . . . . . . . . . . . . . 22
2.3.6 The terminal velocity of a falling mass . . . . . . . . . . . 24
2.3.7 Escape velocity . . . . . . . . . . . . . . . . . . . . . . . . 25
2.3.8 The logistic equation . . . . . . . . . . . . . . . . . . . . . 26
v
vi CONTENTS
3 Second-order odes, constant coecients 29
3.1 The principle of superposition . . . . . . . . . . . . . . . . . . . . 30
3.2 Homogeneous odes . . . . . . . . . . . . . . . . . . . . . . . . . . 30
3.2.1 Real, distinct roots . . . . . . . . . . . . . . . . . . . . . . 31
3.2.2 Complex conjugate, distinct roots . . . . . . . . . . . . . 33
3.2.3 Repeated roots . . . . . . . . . . . . . . . . . . . . . . . . 35
3.3 Inhomogeneous odes . . . . . . . . . . . . . . . . . . . . . . . . . 36
3.3.1 First-order linear inhomogeneous odes revisited . . . . . . 39
3.3.2 Resonance . . . . . . . . . . . . . . . . . . . . . . . . . . . 40
3.3.3 Damped resonance . . . . . . . . . . . . . . . . . . . . . . 42
4 The Laplace transform 45
4.1 Denition and properties . . . . . . . . . . . . . . . . . . . . . . . 45
4.2 Solution of initial value problems . . . . . . . . . . . . . . . . . . 49
4.3 Heaviside and Dirac delta functions . . . . . . . . . . . . . . . . . 51
4.3.1 Heaviside function . . . . . . . . . . . . . . . . . . . . . . 52
4.3.2 Dirac delta function . . . . . . . . . . . . . . . . . . . . . 53
4.4 Discontinuous or impulsive terms . . . . . . . . . . . . . . . . . . 54
5 Series solutions 57
5.1 Ordinary points . . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
5.2 Regular singular points: Euler equations . . . . . . . . . . . . . . 61
5.2.1 Real, distinct roots . . . . . . . . . . . . . . . . . . . . . . 63
5.2.2 Complex conjugate roots . . . . . . . . . . . . . . . . . . 63
5.2.3 Repeated roots . . . . . . . . . . . . . . . . . . . . . . . . 63
6 Systems of equations 65
6.1 Determinants and the eigenvalue problem . . . . . . . . . . . . . 65
6.2 Coupled rst-order equations . . . . . . . . . . . . . . . . . . . . 66
6.2.1 Two distinct real eigenvalues . . . . . . . . . . . . . . . . 66
6.2.2 Complex conjugate eigenvalues . . . . . . . . . . . . . . . 70
6.2.3 Repeated eigenvalues with one eigenvector . . . . . . . . . 72
6.3 Normal modes . . . . . . . . . . . . . . . . . . . . . . . . . . . . 75
7 Bifurcation theory 77
7.1 Fixed points and stability . . . . . . . . . . . . . . . . . . . . . . 77
7.2 One-dimensional bifurcations . . . . . . . . . . . . . . . . . . . . 80
7.2.1 Saddle-node bifurcation . . . . . . . . . . . . . . . . . . . 81
7.2.2 Transcritical bifurcation . . . . . . . . . . . . . . . . . . . 82
7.2.3 Supercritical pitchfork bifurcation . . . . . . . . . . . . . 82
7.2.4 Subcritical pitchfork bifurcation . . . . . . . . . . . . . . 83
7.2.5 Application: a mathematical model of a shery . . . . . . 85
7.3 Two-dimensional bifurcations . . . . . . . . . . . . . . . . . . . . 87
7.3.1 Supercritical Hopf bifurcation . . . . . . . . . . . . . . . . 88
7.3.2 Subcritical Hopf bifurcation . . . . . . . . . . . . . . . . . 88
CONTENTS vii
8 Partial dierential equations 91
8.1 Derivation of the diusion equation . . . . . . . . . . . . . . . . . 91
8.2 Derivation of the wave equation . . . . . . . . . . . . . . . . . . . 92
8.3 Fourier Series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 94
8.4 Solution of the diusion equation . . . . . . . . . . . . . . . . . . 96
8.4.1 Homogeneous boundary conditions . . . . . . . . . . . . . 96
8.4.2 Inhomogeneous boundary conditions . . . . . . . . . . . . 100
8.4.3 Pipe with closed ends . . . . . . . . . . . . . . . . . . . . 101
8.5 Solution of the wave equation . . . . . . . . . . . . . . . . . . . . 104
8.5.1 Plucked string . . . . . . . . . . . . . . . . . . . . . . . . 104
8.5.2 Hammered string . . . . . . . . . . . . . . . . . . . . . . . 106
8.5.3 General initial conditions . . . . . . . . . . . . . . . . . . 106
8.6 The Laplace equation . . . . . . . . . . . . . . . . . . . . . . . . 107
8.6.1 Dirichlet problem for a rectangle . . . . . . . . . . . . . . 107
8.6.2 Dirichlet problem for a circle . . . . . . . . . . . . . . . . 108
viii CONTENTS
Chapter 0
A short mathematical
review
0.1 The trigonometric functions
The Pythagorean trigonometric identity is
sin
2
x + cos
2
x = 1,
and the addition theorems are
sin(x +y) = sin(x) cos(y) + cos(x) sin(y),
cos(x +y) = cos(x) cos(y) sin(x) sin(y).
Also, the values of sinx in the rst quadrant can be remembered by the rule of
quarters, with 0
= 0, 30
= /6, 45
= /4, 60
= /3, 90
= /2:
sin0
=
_
0
4
= 0, sin30
=
_
1
4
=
1
2
, sin45
=
_
2
4
=
2
2
,
sin60
=
_
3
4
=
3
2
, sin90
=
_
4
4
= 1.
The following symmetry properties are also useful:
sin(/2 x) = cos x, cos(/2 x) = sinx;
and
sin(x) = sin(x), cos(x) = cos(x).
0.2 The exponential function and the natural
logarithm
The transcendental number e, approximately 2.71828, is dened as
e = lim
n
_
1 +
1
n
_
n
.
1
2 CHAPTER 0. A SHORT MATHEMATICAL REVIEW
The exponential function exp(x) = e
x
and natural logarithm lnx are inverse
functions satisfying
e
ln x
= x, lne
x
= x.
The usual rules of exponents apply:
e
x
e
y
= e
x+y
, e
x
/e
y
= e
xy
, (e
x
)
p
= e
px
.
The corresponding rules for the logarithmic function are
ln(xy) = lnx + lny, ln(x/y) = lnx lny, lnx
p
= p lnx.
0.3 Denition of the derivative
view lecture
The derivative of the function y = f(x), denoted as f
(x) or dy/dx, is dened

as the slope of the tangent line to the curve y = f(x) at the point (x, y). This
slope is obtained by a limit, and is dened as
f
(x) = lim
h0
f(x +h) f(x)
h
.
0.4 Dierentiating a combination of functions
0.4.1 The sum or dierence rule
The derivative of the sum of f(x) and g(x) is
(f +g)
= f
+g
.
Similarly, the derivative of the dierence is
(f g)
= f
.
0.4.2 The product rule
The derivative of the product of f(x) and g(x) is
(fg)
= f
g +fg
,
and should be memorized as the derivative of the rst times the second plus
the rst times the derivative of the second.
0.4.3 The quotient rule
The derivative of the quotient of f(x) and g(x) is
_
f
g
_
=
f
g fg
g
2
,
and should be memorized as the derivative of the top times the bottom minus
the top times the derivative of the bottom over the bottom squared.
0.5. DIFFERENTIATING ELEMENTARY FUNCTIONS 3
0.4.4 The chain rule
The derivative of the composition of f(x) and g(x) is
(f(g(x)))
= f
(g(x)) g
(x),
and should be memorized as the derivative of the outside times the derivative
of the inside.
0.5 Dierentiating elementary functions
0.5.1 The power rule
The derivative of a power of x is given by
d
dx
x
p
= px
p1
.
0.5.2 Trigonometric functions
The derivatives of sinx and cos x are
(sinx)
= cos x, (cos x)
= sinx.
We thus say that the derivative of sine is cosine, and the derivative of cosine
is minus sine. Notice that the second derivatives satisfy
(sinx)
= sinx, (cos x)
= cos x.
0.5.3 Exponential and natural logarithm functions
The derivative of e
x
and lnx are
(e
x
)
= e
x
, (lnx)
=
1
x
.
0.6 Denition of the integral
The denite integral of a function f(x) > 0 from x = a to b (b > a) is dened
as the area bounded by the vertical lines x = a, x = b, the x-axis and the curve
y = f(x). This area under the curve is obtained by a limit. First, the area is
approximated by a sum of rectangle areas. Second, the integral is dened to be
the limit of the rectangle areas as the width of each individual rectangle goes to
zero and the number of rectangles goes to innity. This resulting innite sum
is called a Riemann Sum, and we dene
_
b
a
f(x)dx = lim
N
N
n=1
f(x
n
)x, (1)
where x
n
= a +(n 1)x and x = (b a)/N. The symbols on the left-hand-
side of (1) are read as the integral from a to b of f of x dee x. The Riemann
Sum denition is extended to all values of a and b and for all values of f(x)
(positive and negative). Accordingly,
_
a
b
f(x)dx =
_
b
a
f(x)dx and
_
b
a
(f(x))dx =
_
b
a
f(x)dx.
Also, if a < b < c, then
_
c
a
f(x)dx =
_
b
a
f(x)dx +
_
c
b
f(x)dx,
which states (when f(x) > 0) that the total area equals the sum of its parts.
0.7 The fundamental theorem of calculus
Using the denition of the derivative, we dierentiate the following integral:
d
dx
_
x
a
f(s)ds = lim
h0
_
x+h
a
f(s)ds
_
x
a
f(s)ds
h
= lim
h0
_
x+h
x
f(s)ds
h
= lim
h0
hf(x)
h
= f(x).
This result is called the fundamental theorem of calculus, and provides a con-
nection between dierentiation and integration.
The fundamental theorem teaches us how to integrate functions. Let F(x)
be a function such that F
(x) = f(x). We say that F(x) is an antiderivative of

f(x). Then from the fundamental theorem and the fact that the derivative of a
constant equals zero,
F(x) =
_
x
a
f(s)ds +c.
Now, F(a) = c and F(b) =
_
b
a
f(s)ds + F(a). Therefore, the fundamental
theorem shows us how to integrate a function f(x) provided we can nd its
antiderivative:
_
b
a
f(s)ds = F(b) F(a).
Unfortunately, nding antiderivatives is much harder than nding derivatives,
and indeed, most complicated functions can not be integrated analytically.
0.8 Denite and indenite integrals
The Riemann sum denition of an integral is called a denite integral. It is
convenient to also dene an indenite integral by
_
f(x)dx = F(x),
where F(x) is the antiderivative of f(x).
0.9. INDEFINITE INTEGRALS OF ELEMENTARY FUNCTIONS 5
0.9 Indenite integrals of elementary functions
From our known derivatives of elementary functions, we can determine some
simple indenite integrals. The power rule gives us
_
x
n
dx =
x
n+1
n + 1
+c, n = 1.
When n = 1, and x is positive, we have
_
1
x
dx = lnx +c.
If x is negative, using the chain rule we have
d
dx
ln(x) =
1
x
.
Therefore, since
|x| =
_
x if x < 0;
x if x > 0,
we can generalize our indenite integral to strictly positive or strictly negative
x:
_
1
x
dx = ln|x| +c.
Trigonometric functions can also be integrated:
_
cos xdx = sinx +c,
_
sinxdx = cos x +c.
Easily proved identities are an addition rule:
_
(f(x) +g(x)) dx =
_
f(x)dx +
_
g(x)dx;
and multiplication by a constant:
_
Af(x)dx = A
_
f(x)dx.
This permits integration of functions such as
_
(x
2
+ 7x + 2)dx =
x
3
3
+
7x
2
2
+ 2x +c,
and
_
(5 cos x + sinx)dx = 5 sinx cos x +c.
0.10 Substitution
More complicated functions can be integrated using the chain rule. Since
d
dx
f(g(x)) = f
(g(x)) g
(x),
we have
_
f
(g(x)) g
(x)dx = f(g(x)) +c.

This integration formula is usually implemented by letting y = g(x). Then one
writes dy = g
(x)dx to obtain
_
f
(g(x))g
(x)dx =
_
f
(y)dy
= f(y) +c
= f(g(x)) +c.
0.11 Integration by parts
Another integration technique makes use of the product rule for dierentiation.
Since
(fg)
= f
g +fg
,
we have
f
g = (fg)
fg
.
Therefore,
_
f
(x)g(x)dx = f(x)g(x)
_
f(x)g
(x)dx.
Commonly, the above integral is done by writing
u = g(x) dv = f
(x)dx
du = g
(x)dx v = f(x).
Then, the formula to be memorized is
_
udv = uv
_
vdu.
0.12 Taylor series
A Taylor series of a function f(x) about a point x = a is a power series rep-
resentation of f(x) developed so that all the derivatives of f(x) at a match all
the derivatives of the power series. Without worrying about convergence here,
we have
f(x) = f(a) +f
(a)(x a) +
f
(a)
2!
(x a)
2
+
f
(a)
3!
(x a)
3
+. . . .
Notice that the rst term in the power series matches f(a), all other terms
vanishing, the second term matches f
(a), all other terms vanishing, etc. Com-

monly, the Taylor series is developed with a = 0. We will also make use of the
Taylor series in a slightly dierent form, with x = x
+ and a = x
:
f(x
+) = f(x
) +f
(x
) +
f
(x
)
2!

2
+
f
(x
)
3!

3
+. . . .
0.13. COMPLEX NUMBERS 7
Another way to view this series is that of g() = f(x
+ ), expanded about
= 0.
Taylor series that are commonly used include
e
x
= 1 +x +
x
2
2!
+
x
3
3!
+. . . ,
sinx = x
x
3
3!
+
x
5
5!
. . . ,
cos x = 1
x
2
2!
+
x
4
4!
. . . ,
1
1 +x
= 1 x +x
2
. . . , for |x| < 1,
ln(1 +x) = x
x
2
2
+
x
3
3
. . . , for |x| < 1.
A Taylor series of a function of several variables can also be developed. Here,
all partial derivatives of f(x, y) at (a, b) match all the partial derivatives of the
power series. With the notation
f
x
=
f
x
, f
y
=
f
y
, f
xx
=

2
f
x
2
, f
xy
=

2
f
xy
, f
yy
=

2
f
y
2
, etc.,
we have
f(x, y) = f(a, b) +f
x
(a, b)(x a) +f
y
(a, b)(y b)
+
1
2!
_
f
xx
(a, b)(x a)
2
+ 2f
xy
(a, b)(x a)(y b) +f
yy
(a, b)(y b)
2
_
+. . .
0.13 Complex numbers
view lecture
We dene the imaginary number i to be one of the two numbers that satises
the rule (i)
2
= 1, the other number being i. Formally, we write i =
1. A
complex number z is written as
z = x +iy,
where x and y are real numbers. We call x the real part of z and y the imaginary
part and write
x = Re z, y = Imz.
Two complex numbers are equal if and only if their real and imaginary parts
are equal.
The complex conjugate of z = x +iy, denoted as z, is dened as
z = x iy.
We have
z z = (x +iy)(x iy)
= x
2
+y
2
;
and we dene the absolute value of z, also called the modulus of z, by
|z| = (z z)
1/2
=
_
x
2
+y
2
.
We can add, subtract, multiply and divide complex numbers to get new
complex numbers. With z = x + iy and w = s + it, and x, y, s, t real numbers,
we have
z +w = (x +s) +i(y +t); z w = (x s) +i(y t);
zw = (x +iy)(s +it)
= (xs yt) +i(xt +ys);
z
w
=
z w
w w
=
(x +iy)(s it)
s
2
+t
2
=
(xs +yt)
s
2
+t
2
+i
(ys xt)
s
2
+t
2
.
Furthermore,
|zw| =
_
(xs yt)
2
+ (xt +ys)
2
=
_
(x
2
+y
2
)(s
2
+t
2
)
= |z||w|;
and
zw = (xs yt) i(xt +ys)
= (x iy)(s it)
= z w.
Similarly
z
w
=
|z|
|w|
, (
z
w
) =
z
w
.
Also, z +w = z + w. However, |z + w| |z| + |w|, a theorem known as the
triangle inequality.
It is especially interesting and useful to consider the exponential function of
an imaginary argument. Using the Taylor series expansion of an exponential
function, we have
e
i
= 1 + (i) +
(i)
2
2!
+
(i)
3
3!
+
(i)
4
4!
+
(i)
5
5!
. . .
=
_
1

2
2!
+

4
4!
. . .
_
+i
_

3
3!
+

5
5!
+. . .
_
= cos +i sin.
0.13. COMPLEX NUMBERS 9
Since cos = 1 and sin = 0, we derive the celebrated Eulers identity
e
i
+ 1 = 0,
that links ve fundamental numbers: 0, 1, i, e and using three basic mathe-
matical operations only once: addition, multiplication and exponentiation.
The complex number z can be represented in the complex plane with Re z
as the x-axis and Imz as the y-axis. This leads to the polar representation of
z = x +iy:
z = re
i
,
where r = |z| and tan = y/x. We dene arg z = . Note that is not unique,
though it is conventional to choose the value such that < , and = 0
when r = 0.
Useful trigonometric relations can be derived using e
i
and properties of the
exponential function. The addition law can be derived from
e
i(x+y)
= e
ix
e
iy
.
We have
cos(x +y) + i sin(x +y) = (cos x +i sinx)(cos y +i siny)
= (cos xcos y sinxsiny) +i(sinxcos y + cos xsiny);
yielding
cos(x +y) = cos xcos y sinxsiny, sin(x +y) = sinxcos y + cos xsiny.
De Moivres Theorem derives from e
in
= (e
i
)
n
, yielding the identity
cos(n) +i sin(n) = (cos +i sin)
n
.
For example, if n = 2, we derive
cos 2 +i sin2 = (cos +i sin)
2
= (cos
2
sin
2
) + 2i cos sin.
Therefore,
cos 2 = cos
2
sin
2
, sin2 = 2 cos sin.
With a little more manipulation using cos
2
+ sin
2
= 1, we can derive
cos
2
=
1 + cos 2
2
, sin
2
=
1 cos 2
2
,
which are useful formulas for determining
_
cos
2
d =
1
4
(2 + sin2) +c,
_
sin
2
d =
1
4
(2 sin2) +c,
and
_
2
0
sin
2
d =
_
2
0
cos
2
d = .
Chapter 1
Introduction to odes
1.1 Denition of an ode
A dierential equation is an equation for a function that relates the values of
the function to the values of its derivatives. An ordinary dierential equation
(ode) is a dierential equation for a function of a single variable, e.g., x(t), while
a partial dierential equation (pde) is a dierential equation for a function of
several variables, e.g., v(x, y, z, t). An ode contains ordinary derivatives and a
pde contains partial derivatives. Typically, pdes are much harder to solve than
odes.
1.2 The simplest type of dierential equation
view lecture
The simplest ordinary dierential equations can be integrated directly by nding
antiderivatives. These simplest odes have the form
d
n
x
dt
n
= G(t),
where the derivative of x = x(t) can be of any order, and the right-hand-side
may depend only on the independent variable t. As an example, consider a mass
falling under the inuence of constant gravity, such as approximately found on
the Earths surface. Newtons law, F = ma, results in the equation
m
d
2
x
dt
2
= mg,
where x is the height of the object above the ground, m is the mass of the
object, and g = 9.8 meter/sec
2
is the constant gravitational acceleration. As
Galileo suggested, the mass cancels from the equation, and
d
2
x
dt
2
= g.
Here, the right-hand-side of the ode is a constant. The rst integration, obtained
by antidierentiation, yields
dx
dt
= Agt,
11
12 CHAPTER 1. INTRODUCTION TO ODES
where A is the constant of integration, and the second integration gives
x = B +At
1
2
gt
2
,
with B another constant of integration. The two constants of integration A and
B can be determined from initial conditions. If we know that the initial height
of the mass is x
0
, and the initial velocity is v
0
, then the correct initial conditions
are
x(0) = x
0
,
dx
dt
(0) = v
0
.
Substitution of these initial conditions into the equations for dx/dt and x allows
us to solve for A and B to determine the unique solution that satises both the
ode and the initial conditions:
x(t) = x
0
+v
0
t
1
2
gt
2
. (1.1)
As a real example, suppose we drop a ball o the top of a 50 meter building.
How long will it take the ball to hit the ground? This question requires solution
of (1.1) for the time T it takes for x(T) = 0, given x
0
= 50 meter and v
0
= 0.
Solving for T:
T =
_
2x
0
g
=
_
2 50
9.8
sec
3.2sec.
Chapter 2
First-order dierential
equations
Reference: Boyce and DiPrima, Chapter 2
The general rst-order dierential equation for the function y = y(x) is written
as
dy
dx
= f(x, y), (2.1)
where f(x, y) can be any function of the independent variable x and the de-
pendent variable y. It is not always possible to nd an analytical solution for
y = y(x). Provided f(x, y) is a well-behaved function, it is always possible
to determine a unique numerical solution given an initial value y(x
0
) = y
0
.
The dierential equation (2.1) gives us the slope f(x
0
, y
0
) of the tangent line
to the solution curve y = y(x) at the point (x
0
, y
0
). With a small step size
x, the initial condition (x
0
, y
0
) can be marched forward in the x-coordinate
to x = x
0
+ x, and along the tangent line using Eulers method to obtain the
y-coordinate
y(x
0
+ x) = y(x
0
) + xf(x
0
, y
0
).
This solution (x
0
+x, y
0
+y) then becomes the new initial condition and is
marched forward in the x-coordinate another x, and along the newly deter-
mined tangent line. For small enough x, the numerical solution converges to
the exact solution.
Analytical solution is possible, however, for special forms of (2.1). Here,
we will learn techniques for solving analytically separable and linear rst-order
equations.
2.1 Separable equations
view lecture
A rst-order ode is separable if it can be written in the form
g(y)
dy
dx
= f(x), y(x
0
) = y
0
, (2.2)
13
14 CHAPTER 2. FIRST-ORDER ODES
where the function g(y) is independent of x and f(x) is independent of y. Inte-
gration from x
0
to x results in
_
x
x
0
g(y(x))y
(x)dx =
_
x
x
0
f(x)dx.
The integral on the left can be transformed by substituting u = y(x), du =
y
(x)dx, and changing the lower and upper limits of integration to y(x
0
) = y
0
and y(x) = y. Therefore,
_
y
y
0
g(u)du =
_
x
x
0
f(x)dx,
and since u is a dummy variable of integration, we can write this in the equivalent
form
_
y
y
0
g(y)dy =
_
x
x
0
f(x)dx. (2.3)
A simpler procedure that also yields (2.3 is to treat dy/dx in (2.2) like a fraction.
Multiplying (2.2) by dx results in
g(y)dy = f(x)dx,
which is a separated equation with all the dependent variables on the left-side,
and all the independent variables on the right-side. Equation (2.3) then results
directly upon integration.
Example: Solve
dy
dx
+
1
2
y =
3
2
, with y(0) = 2.
We rst manipulate the dierential equation to the form
dy
dx
=
1
2
(3 y), (2.4)
and then treat dy/dx as if it was a fraction to separate variables:
dy
3 y
=
1
2
dx.
We integrate the right-side from the initial condition x = 0 to x and the left-side
from the initial condition y(0) = 2 to y. Accordingly,
_
y
2
dy
3 y
=
1
2
_
x
0
dx. (2.5)
The integrals in (2.5) need to be done. Note that y(x) < 3 for nite x or the
integral on the left-side diverges. Therefore, 3 y > 0 and integration yields
ln(3 y) |
y
2
=
1
2
x |
x
0
,
ln(3 y) =
1
2
x,
3 y = e
1
2
x
,
y = 3 e
1
2
x
.
2.1. SEPARABLE EQUATIONS 15
0 1 2 3 4 5 6 7
0
1
2
3
4
5
6
x
y
dy/dx + y/2 = 3/2
Figure 2.1: Solution of the following ode:
dy
dx
+
1
2
y =
3
2
.
Since this is our rst nontrivial analytical solution, it is prudent to check our
result. We do this by dierentiating our solution:
dy
dx
=
1
2
e
1
2
x
=
1
2
(3 y);
and checking the initial conditions, y(0) = 3 e
0
= 2. Therefore, our solution
satises both the original ode and the initial condition.
Example: Solve
dy
dx
+
1
2
y =
3
2
, with y(0) = 4.
This is the identical dierential equation as before, but with dierent initial
conditions. We will jump directly to the integration step:
_
y
4
dy
3 y
=
1
2
_
x
0
dx.
Now y(x) > 3, so that y 3 > 0 and integration yields
ln(y 3) |
y
4
=
1
2
x |
x
0
,
ln(y 3) =
1
2
x,
y 3 = e
1
2
x
,
y = 3 +e
1
2
x
.
The solution curves for a range of initial conditions is presented in Fig. 2.1.
All solutions have a horizontal asymptote at y = 3 at which dy/dx = 0. For
y(0) = y
0
, the general solution can be shown to be y(x) = 3+(y
0
3) exp(x/2).
Example: Solve
dy
dx
=
2 cos 2x
3+2y
, with y(0) = 1. (i) For what values of
x > 0 does the solution exist? (ii) For what value of x > 0 is y(x)
maximum?
Notice that the solution of the ode may not exist when y = 3/2, since dy/dx
. We separate variables and integrate from initial conditions:
(3 + 2y)dy = 2 cos 2xdx
_
y
1
(3 + 2y)dy = 2
_
x
0
cos 2xdx
3y +y
2
|
y
1
= sin2x |
x
0
y
2
+ 3y + 2 sin2x = 0
y
=
1
2
[3
1 + 4 sin2x].
Solving the quadratic equation for y has introduced a spurious solution that
does not satisfy the initial conditions. We test:
y
(0) =
1
2
[3 1] =
_
-1;
-2.
Only the + root satises the initial condition, so that the only solution to the
ode and initial conditions is
y =
1
2
[3 +
1 + 4 sin2x]. (2.6)
To determine (i) the values of x > 0 for which the solution exists, we require
1 + 4 sin2x 0,
or
sin2x
1
4
. (2.7)
To evaluate this expression for x, one needs to recall the denition of arcsin, or
sin
1
as found on a typical scientic calculator. The inverse of the function
f(x) = sinx, /2 x /2
is denoted by arcsin. The rst solution with x > 0 of the equation sin2x = 1/4
places 2x in the interval (, 3/2), so to invert this equation using the arcsine
we need to rewrite it as sin( 2x) = 1/4. The solution of this equation is
2x = arcsin
1
4
x =
1
2
_
+ arcsin
1
4
_
.
Therefore the solution exists for 0 x
1
2
_
+ arcsin
1
4
_
= 1.6971 . . ., where
we have used a calculator value for the arcsin(0.25).
To determine (ii) the value of x at which y = y(x) is maximum, we examine
(2.6) directly. The value of y will be maximum when sin2x takes its maximum
value over the interval where the solution exists. This will be when 2x = /2,
or x = /4 = 0.7854 . . ..
The graph of y = y(x) is shown in Fig. 2.2.
2.2. LINEAR EQUATIONS 17
Figure 2.2: Solution of the following ode:
dy
dx
=
2 cos 2x
3+2y
, y(0) = 1.
2.2 Linear equations
view lecture
The rst-order linear dierential equation (linear in y and its derivative) can be
written in the form
dy
dx
+p(x)y = g(x), (2.8)
with the initial condition y(x
0
) = y
0
. Linear rst-order equations can be inte-
grated using an integrating factor (x). First, we multiply (2.8) by (x):
(x)
_
dy
dx
+p(x)y
_
= (x)g(x). (2.9)
Next, we try to determine an integrating factor so that
(x)
_
dy
dx
+p(x)y
_
=
d
dx
[(x)y], (2.10)
and (2.9) becomes
d
dx
[(x)y] = (x)g(x). (2.11)
Equation (2.11) is easily integrated using (x
0
) =
0
and y(x
0
) = y
0
:
(x)y
0
y
0
=
_
x
x
0
(x)g(x)dx,
y =
1
(x)
_
0
y
0
+
_
x
x
0
(x)g(x)dx
_
. (2.12)
It remains to determine (x) from (2.10). Dierentiating and expanding (2.10)
yields
dy
dx
+py =
d
dx
y +
dy
dx
;
and upon simplifying,
d
dx
= p. (2.13)
Equation (2.13) is separable and can be integrated:
_

0
d
=
_
x
x
0
p(x)dx,
ln

0
=
_
x
x
0
p(x)dx,
(x) =
0
exp
__
x
x
0
p(x)
_
dx.
Notice that since
0
cancels out of (2.12), it is customary to assign
0
= 1. The
solution to (2.8) satisfying the initial condition y(x
0
) = y
0
is then traditionally
written as
y =
1
(x)
_
y
0
+
_
x
x
0
(x)g(x)dx
_
,
where
(x) = exp
__
x
x
0
p(x)dx
_
.
This important result nds frequent use in applied mathematics.
Example: Solve
dy
dx
+ 2y = e
x
, with y(0) = 3/4.
Note that this equation is not separable. With p(x) = 2 and g(x) = e
x
, we
have
(x) = exp
__
x
0
2dx
_
= e
2x
,
and
y = e
2x
_
3
4
+
_
x
0
e
2x
e
x
dx
_
= e
2x
_
3
4
+
_
x
0
e
x
dx
_
= e
2x
_
3
4
+ (e
x
1)
_
= e
2x
_
e
x
1
4
_
= e
x
_
1
1
4
e
x
_
.
2.3. APPLICATIONS 19
Example: Solve
dy
dx
2xy = x, with y(0) = 0.
This equation is separable, and we solve it in two ways. First, using an inte-
grating factor with p(x) = 2x and g(x) = x:
(x) = exp
_
2
_
x
0
xdx
_
= e
x
2
,
and
y = e
x
2
_
x
0
xe
x
2
dx.
The integral can be done by substitution with u = x
2
, du = 2xdx:
_
xe
x
2
dx =
1
2
_
e
u
du
=
1
2
e
u
+c
=
1
2
e
x
2
+c,
so that
y =
1
2
e
x
2
_
e
x
2
1
_
=
1
2
_
e
x
2
1
_
.
Second, we integrate by separating variables:
dy
dx
2xy = x,
dy
dx
= x(1 + 2y),
_
y
0
dy
1 + 2y
=
_
x
0
xdx,
1
2
ln(1 + 2y) =
1
2
x
2
,
1 + 2y = e
x
2
,
y =
1
2
_
e
x
2
1
_
.
The results from the two dierent solution methods are the same, and the choice
of method is a personal preference.
2.3 Applications
2.3.1 Compound interest
view lecture
The equation for the growth of an investment with continuous compounding
of interest is a rst-order dierential equation. Let S(t) be the value of the
investment at time t, and let r be the annual interest rate compounded after
every time interval t. The value of the investment at the time t + t is
S(t + t) = S(t) +rtS(t), (2.14)
where rtS(t) is the amount of interest earned during the time period t.
Typical values of these variables are r = 6% per year = 0.06 year
1
and t = 1
month = 1/12 year; in this case rt = 0.005. Rearranging the terms of (2.14),
we obtain
S(t + t) S(t)
t
= rS(t).
To obtain the equation for continuous compounding of interest, we consider the
limit t 0. The resulting dierential equation is
dS
dt
= rS. (2.15)
If S(0) = S
0
is the initial investment, then the solution to (2.15) is determined
by separating variables:
_
S
S
0
dS
S
= r
_
t
0
dt,
ln(S/S
0
) = rt,
S(t) = S
0
e
rt
;
and compounding results in the exponential growth of the initial investment.
The dierential equation (2.15) can be modied to include continuous deposits
or withdrawals, and we will return to this after a brief digression.
2.3.2 Annual percentage yield (APY)
view lecture
The annual percentage yield (APY) expresses an annual rate of interest, taking
into the eect of compounding. In other words, the APY is the equivalent in-
terest rate one would earn if interest was paid only once per year. For example,
to determine the APY of a 6% annual interest rate compounded continuously,
we equate the value of the investment after one year assuming a single inter-
est payment with APY R, or continuous compounding with interest rate r.
Accordingly, we solve
S
0
(1 +R) = S
0
e
r
for R:
R = e
r
1.
For example, if r = 6.00%, then R = 6.18%. The APY is always larger than a
compounded interest rate since compounding earns interest on interest.
2.3.3 The rule of 72
A useful estimate of an investments doubling time is the Rule of 72. For
example, if the APY of an investment is 6%, then the Rule of 72 estimates the
doubling time of that investment to be 72/6 = 12 years, i.e., the principle will
double in approximately 12 years at a 6% APY. To derive the Rule of 72, let
R be the APY of an investment. Then with annual compounding based on the
APY, the investment grows as
S(t) = S
0
(1 +R)
t
,
where t is time, measured in years. The doubling time is given by
S() = 2S
0
= S
0
(1 +R)
,
so that and R satisfy the equation
(1 +R)
= 2.
Apparently, the doubling time depends on the APY R, but commonly R << 1
(e.g., R = 0.06 for an APY of 6%). We take the logarithm of both sides, and
Taylor series expand in R:
ln(1 +R) = ln2,
(R
1
2
R
2
+. . .) = ln2,
=
ln2
R(1
1
2
R +. . .)
,
=
(ln2)(1 +
1
2
R +. . .)
R
,
=
69.31 (1 +
1
2
R +. . .)
100 R
.
For very small R, the rule of 69 would be more appropriate. However, the rule
of 72 is an excellent approximation when
72 = 69.31 (1 +
1
2
R),
or R = 7.76%, an interest rate not far from rates usually considered. Typically,
when mentally calculating, 72 is preferred because of its large number of divisors
spanning a range of APYs: 2, 4, 6, 8, 12, 18, 36.
2.3.4 Compound interest with deposits or withdrawals
We extend the continuous compounding model to include deposits (or with-
drawals). Let k be the annual deposit amount, and suppose that an installment
is deposited after every time interval t. With compounding also occurring
with period t, we have
S(t + t) = S(t) + (rt)S(t) +kt,
where kt is the amount of money deposited (k > 0) or withdrawn (k < 0) at
time t+t. Dierencing and taking the limit t 0, we obtain the dierential
equation for both continuous compounding of interest and continuous rates of
deposits (or withdrawals):
dS
dt
= rS +k,
which can solved with the initial condition S(0) = S
0
, where S
0
is the initial
capital. We can solve either by separating variables or by using an integrating
factor; I solve here by separating variables. Integrating from t = 0 to a nal
time t,
_
S
S
0
dS
rS +k
=
_
t
0
dt,
1
r
ln
_
rS +k
rS
0
+k
_
= t,
rS +k = (rS
0
+k)e
rt
,
S =
rS
0
e
rt
+ke
rt
k
r
,
S = S
0
e
rt
+
k
r
e
rt
_
1 e
rt
_
, (2.16)
where in (2.16), the rst term comes from the initial invested capital and the
second term comes from the deposits (or withdrawals). As a practical example,
we can analyze a simple retirement plan. It is easiest to assume that all amounts
and returns are in real dollars (adjusted for ination). Suppose a 25 year-old
plans to set aside a xed amount every year of his/her working life, invests at
a real return of 6%, and retires at age 65. How much must he/she invest each
year to have $8 000 000 at retirement? We have t = 40 years, S(t) = $8 000 000,
S
0
= 0, r = 0.06 per year. We need to solve (2.16) for k. With S
0
= 0,
k =
rS(t)
e
rt
1
,
k =
0.06 8 000 000
e
0.0640
1
,
= $47 889 year
1
.
To have saved approximately one million US$ at retirement, the worker would
need to save about HK$50 000 per year over his/her working life. Note that
the amount saved over the workers life is 40 $47 899 = $1 915 960, while the
amount earned on the investment (at the assumed 6% real return) is $8 000 000
$1 915 960 = $6 084 040. The amount earned from the investment is about 3
the amount saved, even with the modest real return of 6%. Sound investment
planning is well worth the eort.
2.3.5 Chemical reactions
view lecture
Suppose that two chemicals A and B react to form a product C, which we write
as
A+B
k
C,
where k is called the rate constant of the reaction. For simplicity, we will use
the same symbol C, say, to refer to both the chemical C and its concentration.
The law of mass action says that dC/dt is proportional to the product of the
concentrations A and B, with proportionality constant k; that is,
dC
dt
= kAB. (2.17)
Similarly, the law of mass action enables us to write equations for the time-
derivatives of the reactant concentrations A and B:
dA
dt
= kAB,
dB
dt
= kAB. (2.18)
The ode given by (2.17) can be solved analytically using conservation laws.
From (2.17) and (2.18),
d
dt
(A+C) = 0 = A+C = A
0
,
d
dt
(B +C) = 0 = B +C = B
0
,
where A
0
and B
0
are the initial concentrations of the reactants, and we assume
that no product is present initially. Using the conservation laws, (2.17) becomes
dC
dt
= k(A
0
C)(B
0
C), with C(0) = 0,
which is a nonlinear equation that may be integrated by separating variables.
Separating and integrating, we obtain
_
C
0
dC
(A
0
C)(B
0
C)
= k
_
t
0
= kt. (2.19)
The remaining integral can be done using the method of partial fractions. We
write
1
(A
0
C)(B
0
C)
=
a
A
0
C
+
b
B
0
C
.
Multiplying both sides by (A
0
C)(B
0
C) yields
1 = a(B
0
C) +b(A
0
C)
= (B
0
a +A
0
b) (a +b)C.
Equating coecients of C
0
and C
1
results in the two linear equations
B
0
a +A
0
b = 1, a +b = 0,
which are easily solved for a and b:
a =
1
B
0
A
0
, b =
1
B
0
A
0
.
Therefore,
1
(A
0
C)(B
0
C)
=
1
B
0
A
0
_
1
A
0
C

1
B
0
C
_
,
and the integral becomes (using C < A
0
, B
0
)
_
C
0
dC
(A
0
C)(B
0
C)
=
1
B
0
A
0
_
_
C
0
dC
A
0
C

_
C
0
dC
B
0
C
_
=
1
B
0
A
0
_
ln
_
A
0
C
A
0
_
+ ln
_
B
0
C
B
0
__
=
1
B
0
A
0
ln
_
A
0
(B
0
C)
B
0
(A
0
C)
_
.
Using this integral in (2.19), multiplying by (B
0
A
0
) and exponentiating, we
obtain
A
0
(B
0
C)
B
0
(A
0
C)
= e
(B
0
A
0
)kt
.
Solving for C, we nally obtain
C(t) = A
0
B
0
e
(B
0
A
0
)kt
1
B
0
e
(B
0
A
0
)kt
A
0
,
which appears to be a complicated expression, but has the simple limits
lim
t
C(t) =
_
A
0
if A
0
< B
0
,
B
0
if B
0
< A
0
,
= min(A
0
, B
0
).
Hence, the reaction stops after one of the reactants is depleted; and the nal
concentration of product is equal to the initial concentration of the depleted
reactant.
2.3.6 The terminal velocity of a falling mass
view lecture
Using Newtons law, we model a mass falling under gravity with air resistance.
We assume that the force of air resistance is proportional to the speed of the
mass and opposes the direction of motion. The total force on the mass is
therefore given by F = mg kv, where v is the unidirectional velocity of the
mass. With F = ma and a = dv/dt, we obtain the dierential equation
m
dv
dt
= mg kv. (2.20)
The terminal velocity v
t
of the mass is dened as the asymptotic velocity after
air resistance balances the gravitational force. When the mass is at terminal
velocity, dv/dt = 0 so that
v
t
=
mg
k
.
The approach to the terminal velocity of a mass initially at rest is obtained by
solving (2.20) with initial condition v(0) = 0. The equation is both linear and
separable, and I solve by separating variables:
m
_
v
0
dv
mg +kv
=
_
t
0
dt,
m
k
ln
_
mg +kv
mg
_
= t,
1 +
kv
mg
= e
kt/m
,
v =
mg
k
_
1 e
kt/m
_
.
Therefore, v = v
t
_
1 e
kt/m
_
, and v approaches v
t
as the exponential term
decays to zero.
As an example, a skydiver with the parachute closed may have a terminal ve-
locity of 200km/hr. Assuming m = 100kg and g = 9.8m/s
2
= 127, 008km/hr
2
,
one obtains k = 63, 504kg/hr. One-half of the terminal velocity for free-fall
(100km/hr) is therefore attained when (1 e
kt/m
) = 1/2, or t = mln2/k
4 sec. Approximately 95% of the terminal velocity is attained after 17 sec.
2.3.7 Escape velocity
view lecture
An interesting physical problem is to nd the smallest initial velocity for a
mass on the Earths surface to escape from the Earths gravitational eld, the
so-called escape velocity. Newtons law of universal gravitation asserts that the
gravitational force between two massive bodies is proportional to the product of
the two masses and inversely proportional to the square of the distance between
them. For a mass m a position x above the surface of the Earth, the force on
the mass is given by
F = G
Mm
(R +x)
2
,
where M and R are the mass and radius of the Earth and G is the gravitational
constant. The minus sign occurs because the force on the mass m points in the
opposite direction of increasing x. The approximately constant acceleration g
on the Earths surface corresponds to the absolute value of F/m when x = 0:
g =
GM
R
2
,
and g 9.8 m/s
2
. Newtons law F = ma for the mass m is thus given by
d
2
x
dt
2
=
GM
(R +x)
2
=
g
(1 +x/R)
2
. (2.21)
A useful trick allows us to solve this second-order dierential equation as
a rst-order equation. First, note that d
2
x/dt
2
= dv/dt. If we write v(t) =
v(x(t))considering the velocity of the mass m to be a function of its distance
above the Earthwe have using the chain rule
dv
dt
=
dv
dx
dx
dt
= v
dv
dx
,
where we have used v = dx/dt. Therefore, (2.21) becomes the rst-order ode
v
dv
dx
=
g
(1 +x/R)
2
,
which may be solved assuming an initial velocity v(x = 0) = v
0
when the mass
is shot vertically from the Earths surface. Separating variables and integrating,
we obtain
_
v
v
0
vdv = g
_
x
0
dx
(1 +x/R)
2
.
The left integral is
1
2
(v
2
v
2
0
), and the right integral can be performed using
the substitution u = 1 +x/R, du = dx/R:
_
x
0
dx
(1 +x/R)
2
= R
_
1+x/R
1
du
u
2
=
R
u
1+x/R
1
= R
R
2
x +R
=
Rx
x +R
.
Therefore,
1
2
(v
2
v
2
0
) =
gRx
x +R
,
v
2
= v
2
0
2gRx
x +R
.
The escape velocity is dened as the minimum initial velocity v
0
such that the
mass can escape to innity. Therefore, v
0
= v
escape
when v 0 as x .
Taking this limit, we have
v
2
escape
= lim
x
2gRx
x +R
= 2gR.
Therefore, with R 6350 km and g = 9.8 10
3
3600
2
km/hr
2
, v
escape
=
2gR 40 000 km/hr.

2.3.8 The logistic equation
view lecture
Let N = N(t) be the size of a population and r the growth rate. The Malthusian
growth model (Thomas Malthus, 1766-1834), equivalent to a compound interest
model, is given by
dN
dt
= rN.
When the population growth is constrained by limited resources, a heuristic
modication to the Malthusian growth model results in the Verhulst equation,
dN
dt
= rN(1
N
K
),
where K is the carrying capacity of the environment. Nondimensionalization of
the Verhulst equation using = rt and x = N/K leads to the logistic equation,
dx
d
= x(1 x),
where we may assume the initial condition x(0) = x
0
> 0. Separating variables
and integrating
_
x
x
0
dx
x(1 x)
=
_

0
d.
The integral on the left-hand-side can be done using the method of partial
fractions:
1
x(1 x)
=
a
x
+
b
1 x
=
a + (b a)x
x(1 x)
;
and equating the coecients of the numerators proportional to x
0
and x
1
, we
have a = b = 1. Therefore,
_
x
x
0
dx
x(1 x)
=
_
x
x
0
dx
x
+
_
x
x
0
dx
(1 x)
= ln
x
x
0
ln
1 x
1 x
0
= ln
x(1 x
0
)
x
0
(1 x)
= .
Solving for x, we rst exponentiate both sides and then isolate x:
x(1 x
0
)
x
0
(1 x)
= e
,
x(1 x
0
) = x
0
e
xx
0
e
,
x(1 x
0
+x
0
e
) = x
0
e
,
x =
x
0
x
0
+ (1 x
0
)e
. (2.22)
We observe that for x
0
> 0, as , x() 1.
Chapter 3
Second-order linear
dierential equations with
constant coecients
view lecture
The general second-order linear dierential equation with independent variable
t and dependent variable x = x(t) is given by
x +p(t) x +q(t)x = g(t), (3.1)
where we have used the standard physics notation x = dx/dt and x = d
2
x/dt
2
.
A unique solution of (3.1) requires initial values x(t
0
) = x
0
and x(t
0
) = u
0
.
The equation with constant coecientson which we will devote considerable
eortassumes that p(t) and q(t) are constants, independent of time. The
second-order linear ode is said to be homogeneous if g(t) = 0.
In general, (3.1) can not be solved analytically, and we begin by deriving an
algorithm for numerical solution. We write the second-order ode as a pair of
rst-order odes by dening u = x, and writing the rst-order system as
x = u, (3.2)
u = f(t, x, u), (3.3)
where
f(t, x, u) = p(t)u q(t)x +g(t).
The rst ode, (3.2), gives the slope of the tangent line to the curve x = x(t);
the second ode, (3.3), gives the slope of the tangent line to the curve u = u(t).
Beginning at the initial values (x, u) = (x
0
, u
0
) at the time t = t
0
, we move
along the tangent lines to determine x
1
= x(t
0
+ t) and u
1
= u(t
0
+ t):
x
1
= x
0
+ tu
0
,
u
1
= u
0
+ tf(t
0
, x
0
, u
0
).
The values x
1
and u
1
at the time t
1
= t
0
+t are then used as new initial values
to march the solution forward to time t
2
= t
1
+ t. As long as f(t, x, u) is a
well-behaved function, the numerical solution converges to the unique solution
of the ode as t 0.
29
30 CHAPTER 3. SECOND-ORDER ODES, CONSTANT COEFFICIENTS
3.1 The principle of superposition
view lecture
Consider the second-order linear homogeneous ode:
x +p(t) x +q(t)x = 0; (3.4)
and suppose that X
1
(t) and X
2
(t) are solutions to (3.4). We consider a linear
superposition of X
1
and X
2
by letting
X(t) = c
1
X
1
(t) +c
2
X
2
(t).
The principle of superposition states that X(t) is also a solution of (3.4). To
prove this, we compute
X +p

X +qX = c
1

X
1
+c
2

X
2
+p
_
c
1

X
1
+c
2

X
2
_
+q (c
1
X
1
+c
2
X
2
)
= c
1
_
X
1
+p

X
1
+qX
1
_
+c
2
_
X
2
+p

X
2
+qX
2
_
= c
1
0 +c
2
0
= 0,
since X
1
and X
2
were assumed to be solutions of (3.4). We have therefore shown
that any linear combination of solutions to the second-order linear homogeneous
ode is also a solution.
3.2 Second-order linear homogeneous ode with
constant coecients
view lecture
We now study solutions of the equation
a x +b x +cx = 0, (3.5)
with a, b, and c constants. Such an equation arises for the charge on a capacitor
in an unpowered RLC electrical circuit, or for the position of a freely-oscillating
frictional mass on a spring. The principle of superposition states that any lin-
ear combination of solutions to (3.5) is also a solution. Our solution method,
therefore, is to nd all linearly independent solutions to (3.5), multiply these
solutions by constants and add them. (Proof that two or more solutions are lin-
early independent proceeds by computing a determinant called the Wronskian.
Since this is mainly of theoretical interest and rarely done in practice, I will
omit a discussion of the Wronskian here. Interested students can consult the
textbook or Wikipedia.) Since two initial conditions determine a unique solu-
tion, we expect to nd two linearly independent solutions and for the general
solution of (3.5) to contain two free constants.
Because of the dierential properties of the exponential function, a natural
ansatz, or guess, for the form of the solution to (3.5) is x = e
rt
, where r is a
constant to be determined. Successive dierentiation results in x = re
rt
and
x = r
2
e
rt
, and substitution into (3.5) yields
ar
2
e
rt
+bre
rt
+ce
rt
= 0. (3.6)
3.2. HOMOGENEOUS ODES 31
Our choice of exponential function is now rewarded by the explicit cancellation
of e
rt
in (3.6). The result is a quadratic equation for the unknown constant r:
ar
2
+br +c = 0. (3.7)
Our ansatz has thus converted a dierential equation into an algebraic equation.
Equation (3.7) is called the characteristic equation of (3.5). Using the quadratic
formula, the two solutions to the characteristic equation (3.7) are given by
r
=
1
2a
_
b
_
b
2
4ac
_
.
There are three cases to consider: (1) if b
2
4ac > 0, then the two roots are real
and distinct; (2) if b
2
4ac < 0, then the two roots are distinct and complex
conjugates of each other; (3) if b
2
4ac = 0, then the two roots are degenerate
and there is only one real root. We will consider these three cases in turn.
3.2.1 Real, distinct roots
When r
+
= r
are real roots, then the general solution to (3.5) can be written
as a linear superposition of the two solutions e
r
+
t
and e
r
t
; that is,
x(t) = c
1
e
r
+
t
+c
2
e
r
t
.
The unknown constants c
1
and c
2
can then be determined by the initial condi-
tions x(t
0
) = x
0
and x(t
0
) = u
0
. We now present two examples.
Example 1: Solve x + 5 x + 6x = 0 with x(0) = 2, x(0) = 3, and nd the
maximum value attained by x.
We take as our ansatz x = e
rt
and obtain the characteristic equation
r
2
+ 5r + 6 = 0,
which factors to
(r + 3)(r + 2) = 0.
The general solution to the ode is thus
x(t) = c
1
e
2t
+c
2
e
3t
.
The solution for x obtained by dierentiation is
x(t) = 2c
1
e
2t
3c
2
e
3t
.
Use of the initial conditions then results in two equations for the two unknown
constant c
1
and c
2
:
c
1
+c
2
= 2,
2c
1
3c
2
= 3.
Adding three times the rst equation to the second equation yields c
1
= 9; and
the rst equation then yields c
2
= 2 c
1
= 7. Therefore, the unique solution
that satises both the ode and the initial conditions is
x(t) = 9e
2t
7e
3t
.
Note that although both exponential terms decay in time, their sum increases
initially since x(0) > 0. The maximum value of x occurs at the time t
m
when
x = 0, or
18e
2t
m
+ 21e
3t
m
= 0,
which may be solved for t
m
:
t
m
= ln
_
7
6
_
.
The maximum x
m
= x(t
m
) is then determined to be
x
m
= 9
_
6
7
_
2
7
_
6
7
_
3
=
_
6
7
_
2
(9 6)
=
108
49
.
Example 2: Solve x x = 0 with x(0) = x
0
, x(0) = u
0
.
view lecture
Again our ansatz is x = e
rt
, and we obtain the characteristic equation
r
2
1 = 0,
so that r
= 1. Therefore, the general solution for x is

x(t) = c
1
e
t
+c
2
e
t
,
and the derivative satises
x(t) = c
1
e
t
c
2
e
t
.
Initial conditions are satised when
c
1
+c
2
= x
0
,
c
1
c
2
= u
0
.
Adding and subtracting these equations, we determine
c
1
=
1
2
(x
0
+u
0
) , c
2
=
1
2
(x
0
u
0
) ,
so that after rearranging terms
x(t) = x
0
_
e
t
+e
t
2
_
+u
0
_
e
t
e
t
2
_
.
The terms in parentheses are the usual denitions of the hyperbolic cosine and
sine functions; that is,
cosht =
_
e
t
+e
t
2
_
, sinht =
_
e
t
e
t
2
_
.
Our solution can therefore be rewritten as
x(t) = x
0
cosht +u
0
sinht.
Note that the hyperbolic trigonometric functions satisfy the dierential equa-
tions
d
dt
sinht = cosht,
d
dt
cosht = sinht.
3.2.2 Complex conjugate, distinct roots
view lecture
We now consider a characteristic equation (3.7) with b
2
4ac < 0, so the roots
occur as complex conjugate pairs. With
=
b
2a
, =
1
2a
_
4ac b
2
,
the two roots of the characteristic equation are + i and i. Therefore,
the general solution to the ode is given by
x(t) = c
1
e
(+i)t
+c
2
e
(i)t
= e
t
_
c
1
e
it
+c
2
e
it
_
. (3.8)
We know, however, that the solution to the dierential equation (3.5) must be
real. Therefore, c
1
and c
2
in (3.8) must be complex and their values are such
that x(t) is real, or x = x. Therefore, we have
c
1
e
it
+c
2
e
it
= c
1
e
it
+c
2
e
it
,
or rewriting,
(c
1
c
2
)e
it
= (c
1
c
2
)e
it
= (c
1
c
2
)e
it
.
Therefore, (c
1
c
2
)e
it
is real for all t, which can only be true if
c
1
= c
2
;
that is, c
1
and c
2
are complex conjugate pairs. Therefore,
x(t) = e
t
_
c
1
e
it
+c
1
e
it
_
= 2e
t
Re
_
c
1
e
it
_
,
where we have used the identity z + z = 2Re(z). Now,
c
1
= Re(c
1
) +iIm(c
1
), e
it
= cos t +i sint,
so that
Re
_
c
1
e
it
_
= Re(c
1
) cos t Im(c
1
) sint,
and
x(t) = e
t
(2Re(c
1
) cos t 2Im(c
1
) sint) . (3.9)
We may rewrite (3.9) by dening new real constants A and B such that
A = 2Re(c
1
), B = 2Im(c
1
).
The solution for x(t) can then be rewritten as
x(t) = e
t
(Acos t +Bsint) . (3.10)
It is best to memorize this result. The real part of the roots of the characteristic
equation goes into the exponential function; the imaginary part goes into the
argument of cosine and sine.
Example 1: Solve x +x = 0 with x(0) = x
0
and x(0) = u
0
.
The characteristic equation is
r
2
+ 1 = 0,
with roots
r
= i.
The general solution of the ode is therefore
x(t) = Acos t +Bsint.
The derivative is
x(t) = Asint +Bcos t.
Applying the initial conditions:
x(0) = A = x
0
, x(0) = B = u
0
;
so that the nal solution is
x(t) = x
0
cos t +u
0
sint.
Example 2: Solve x + x +x = 0 with x(0) = 1 and x(0) = 0.
r
2
+r + 1 = 0,
with roots
r
=
1
2
i
3
2
.
x(t) = e
1
2
t
_
Acos
3
2
t +Bsin
3
2
t
_
.
The derivative is
x(t) =
1
2
e
1
2
t
_
Acos
3
2
t +Bsin
3
2
t
_
+
3
2
e
1
2
t
_
Asin
3
2
t +Bcos
3
2
t
_
.
Applying the initial conditions x(0) = 1 and x(0) = 0:
A = 1,
1
2
A+
3
2
B = 0;
or
A = 1, B =
3
3
.
Therefore,
x(t) = e
1
2
t
_
cos
3
2
t +
3
3
sin
3
2
t
_
.
3.2.3 Repeated roots
view lecture (method of reduction of order)
Finally, we consider a characteristic equation (3.7) with only one real root (b
2
4ac = 0), given by

r =
b
2a
,
yielding only a single solution to the ode:
x
1
(t) = exp
_
bt
2a
_
. (3.11)
To satisfy two initial conditions, a second independent solution must be found.
We can determine this solution by considering the solution for x(t) given by
(3.10) when the roots of the characteristic equation are i, and taking the
limit 0. First, applying initial conditions x(0) = x
0
and x(0) = u
0
to (3.10)
to solve for A and B, we nd
x(t) = e
t
_
x
0
cos t +
u
0
x
0
sint
_
.
Taking the limit as 0,
lim
0
x(t) = e
t
(x
0
+ (u
0
x
0
)t) .
The second solution can be seen to be a constant, u
0
x
0
, times t times the
rst solution, e
t
. Our general solution to the ode (3.5) when b
2
4ac = 0 is
therefore
x(t) = (c
1
+c
2
t) exp
_
bt
2a
_
.
The main result to be remembered is that for the case of repeated roots, the
second solution is t times the rst solution.
Example: Solve x + 2 x +x = 0 with x(0) = 1 and x(0) = 0.
view lecture
r
2
+ 2r + 1 = (r + 1)
2
= 0,
which has a repeated root given by r = 1. Therefore, the general solution to
the ode is
x(t) = c
1
e
t
+c
2
te
t
,
with derivative
x(t) = c
1
e
t
+c
2
e
t
c
2
te
t
.
Applying the initial conditions:
c
1
= 1,
c
1
+c
2
= 0,
so that c
1
= c
2
= 1. Therefore, the solution is
x(t) = (1 +t)e
t
.
3.3 Second-order linear inhomogeneous ode
view lecture
We now consider the general second-order linear inhomogeneous ode (3.1):
x +p(t) x +q(t)x = g(t), (3.12)
with initial conditions x(t
0
) = x
0
and x(t
0
) = u
0
. There is a three-step solution
method when the inhomogeneous term g(t) = 0. (i) Find the general solution
of the homogeneous equation
x +p(t) x +q(t)x = 0. (3.13)
Let us denote the homogeneous solution by
x
h
(t) = c
1
x
1
(t) +c
2
x
2
(t),
where x
1
and x
2
are linearly independent solutions of (3.13), and c
1
and c
2
are as yet undetermined constants. (ii) Find any particular solution x
p
of the
inhomogeneous equation (3.12). This particular solution can be found, for ex-
ample, by the the method of undetermined coecients discussed below, when
g(t) is a combination of polynomials, exponentials, sines and cosines. (iii) Write
the general solution of (3.12) as the sum of the homogeneous and particular
solutions,
x(t) = x
h
(t) +x
p
(t), (3.14)
and apply the initial conditions to determine the constants c
1
and c
2
. Note that
because of the linearity of (3.12),
x +p x +qx =
d
2
dt
2
(x
h
+x
p
) +p
d
dt
(x
h
+x
p
) +q(x
h
+x
p
)
= ( x
h
+p x
h
+qx
h
) + ( x
p
+p x
p
+qx
p
)
= 0 +g
= g,
so that (3.14) solves (3.12), and the two free constants in x
h
are determined by
the initial conditions. We now illustrate this solution method by an example.
3.3. INHOMOGENEOUS ODES 37
Example: Solve x 3 x 4x = 3e
2t
with x(0) = 1 and x(0) = 0.
First, we solve the homogeneous equation. The characteristic equation is
r
2
3r 4 = (r 4)(r + 1)
= 0,
so that
x
h
(t) = c
1
e
4t
+c
2
e
t
.
Second, we nd a particular solution of the inhomogeneous equation. The
method of undetermined coecients uses an ansatz of the form
x
p
(t) = Ae
2t
, (3.15)
where A is the so-called undetermined coecient. The form of the particular
solution is chosen such that the exponential will cancel out of both sides of the
ode. Upon substitution of x
p
into the ode, dierentiating using the chain rule
and canceling the exponential, we obtain:
4A6A4A = 3,
from which we determine A = 1/2. Obtaining a solution for A independent of
t justies the ansatz (3.15). Third, we write the general solution to the ode as
the sum of the homogeneous and particular solutions, and determine c
1
and c
2
that satisfy the initial conditions. We have
x(t) = c
1
e
4t
+c
2
e
t
1
2
e
2t
;
and taking the derivative,
x(t) = 4c
1
e
4t
c
2
e
t
e
2t
.
Applying the initial conditions,
c
1
+c
2
1
2
= 1,
4c
1
c
2
1 = 0;
or
c
1
+c
2
=
3
2
,
4c
1
c
2
= 1.
This system of equations can be solved for c
1
by adding the equations to obtain
c
1
= 1/2, after which c
2
can be determined from the rst equation to obtain
c
2
= 1. Therefore, the solution for x(t) that satises both the ode and the initial
conditions is given by
x(t) =
1
2
e
4t
1
2
e
2t
+e
t
=
1
2
e
4t
_
1 e
2t
+ 2e
5t
_
,
where we have grouped the terms in the solution to better display the asymp-
totics for large t. We now discuss several examples using the method of unde-
termined coecients to nd particular solutions.
Example: Find a particular solution of x 3 x 4x = 2 sint.
Although it is possible to solve for a particular solution with the real ansatz
x
p
(t) = Acos t +Bsint,
where the argument of cosine and sine must agree with the argument of sine
in the inhomogeneous term, here, I prefer to make use of the relation e
it
=
cos t +i sint. We introduce the complex solution z(t) by letting
z(t) = x(t) +iy(t),
where x(t) = Re{z(t)} is the real solution to the ode. We then replace the
original real ode to be solved, using sint = Re{ie
it
}, by
z 3 z 4z = 2ie
it
. (3.16)
Notice that (3.16) is actually two real odes, obtained by taking the real and
imaginary parts of both sides:
x 3 x 4x = 2 sint, y 3 y 4y = 2 cos t;
and that we are only interested in the solution of the rst ode. In (3.16), we now
have an exponential as the inhomogeneous term, and we can make the ansatz
z
p
(t) = Ae
it
,
where we now expect A to be complex. Upon substitution into the ode (3.16)
and using i
2
= 1:
A3iA4A = 2i,
or solving for A:
A =
2i
5 + 3i
=
2i(5 3i)
(5 + 3i)(5 3i)
=
6 + 10i
34
=
3 + 5i
17
.
Therefore,
x
p
= Re{z
p
}
= Re
_
1
17
(3 + 5i)(cos t +i sint)
_
=
1
17
(3 cos t 5 sint).
Example: Find a particular solution of x + x 2x = t
2
.
view lecture
The correct ansatz here is the polynomial
x
p
(t) = At
2
+Bt +C.
Upon substitution into the ode
2A+ 2At +B 2At
2
2Bt 2C = t
2
or
2At
2
+ 2(AB)t + (2A+B 2C)t
0
= t
2
.
Equating powers of t:
2A = 1, 2(AB) = 0, 2A+B 2C = 0.
Therefore,
A =
1
2
, B =
1
2
, C =
3
4
;
and our particular solution is
x
p
(t) =
1
2
t
2
1
2
t
3
4
.
3.3.1 First-order linear inhomogeneous odes revisited
view lecture
The rst-order linear ode can be solved by use of an integrating factor. Here I
show that odes having constant coecients can be solved by our newly learned
solution method.
Example: Solve x + 2x = e
t
with x(0) = 3/4.
Rather than using an integrating factor, we follow the three-step approach: (i)
nd the general homogeneous solution; (ii) nd any particular solution; (iii) add
them and satisfy initial conditions. Accordingly, we try the ansatz x
h
(t) = e
rt
for the homogeneous ode x + 2x = 0 and nd
r + 2 = 0, or r = 2.
To nd a particular solution, we try the ansatz x
p
(t) = Ae
t
, and upon substi-
tution
A+ 2A = 1, or A = 1.
Therefore, the general solution to the ode is
x(t) = ce
2t
+e
t
.
The single initial condition determines the unknown constant c:
x(0) =
3
4
= c + 1,
so that c = 1/4. Hence,
x(t) = e
t
1
4
e
2t
= e
t
_
1
1
4
e
t
_
.
3.3.2 Resonance
view lecture
Resonance occurs when the frequency of the inhomogeneous term matches the
frequency of the homogeneous solution. To illustrate resonance in its simplest
embodiment, we consider the second-order linear inhomogeneous ode
x +
2
0
x = f cos t, x(0) = x
0
, x(0) = u
0
. (3.17)
The homogeneous equation has characteristic equation
r
2
+
2
0
= 0,
so that r
= i
0
. Therefore,
x
h
(t) = c
1
cos
0
t +c
2
sin
0
t. (3.18)
To nd a particular solution, we replace the ode (3.17) by its complex equiv-
alent, with x = Re(z):
z +
2
0
z = fe
it
.
Using the ansatz z
p
= Ae
it
, we obtain
2
A+
2
0
A = f,
or
A =
f
2
0
2
,
so that
x
p
=
f
2
0
2
cos t.
Therefore, our general solution is
x(t) = c
1
cos
0
t +c
2
sin
0
t +
f
2
0
2
cos t,
with derivative
x(t) =
0
(c
2
cos
0
t c
1
sin
0
t)
f
2
0
2
sint.
Initial conditions are satised when
x
0
= c
1
+
f
2
0
2
,
u
0
= c
2
0
,
so that
c
1
= x
0
2
0
2
, c
2
=
u
0
0
.
Therefore, the solution to the ode that satises the initial conditions is
x(t) =
_
x
0
2
0
2
_
cos
0
t +
u
0
0
sin
0
t +
f
2
0
2
cos t
= x
0
cos
0
t +
u
0
0
sin
0
t +
f(cos t cos
0
t)
2
0
2
.
Resonance occurs in the limit
0
; that is, the frequency of the inhomo-
geneous term (the external force) matches the frequency of the homogeneous
solution (the internal, natural frequency of oscillation). By LHospitals rule,
the limit of the term proportional to f is found by dierentiating with respect
to :
lim
0
f(cos t cos
0
t)
2
0
2
= lim
0
ft sint
2
=
ft sin
0
t
2
0
. (3.19)
At resonance, the term proportional to the amplitude f of the inhomogenous
term increases linearly with t, resulting in larger-and-larger amplitudes of os-
cillation for x(t). In general, if the inhomogeneous term is a solution of the
homogeneous equation, then the correct ansatz for the particular solution is the
homogeneous solution multiplied by t.
To illustrate this further, we return to the original ode, now assumed to be
exactly at resonance,
x +
2
0
x = f cos
0
t,
and nd a particular solution directly. The particular solution is the real part
of the particular solution of
z +
2
0
z = fe
i
0
t
,
and because the inhomogeneous term is a solution of the homogeneous equation,
we take as our ansatz
z
p
= Ate
i
0
t
.
We have
z
p
= Ae
i
0
t
(1 +i
0
t) , z
p
= Ae
i
0
t
_
2i
0
2
0
t
_
;
and upon substitution into the ode
z
p
+
2
0
z
p
= Ae
i
0
t
_
2i
0
2
0
t
_
+
2
0
Ate
i
0
t
= 2i
0
Ae
i
0
t
= fe
i
0
t
.
Therefore,
A =
f
2i
0
,
and
x
p
= Re{
ft
2i
0
e
i
0
t
}
=
ft sin
0
t
2
0
,
the same result as (3.19).
Example: Find a particular solution of x 3 x 4x = 5e
t
.
The homogeneous equation has solution
x
h
= c
1
e
4t
+c
2
e
t
,
so that the inhomogeneous term is a solution of the homogeneous equation. To
nd a particular solution, we therefore take as our ansatz
x = Ate
t
,
with rst- and second-derivatives given by
x = Ae
t
(1 t), x = Ae
t
(2 +t).
Substitution into the dierential equation yields
Ae
t
(2 +t) 3Ae
t
(1 t) 4Ate
t
= 5e
t
,
which after cancelling terms yields 5A = 5, or A = 1. Therefore, the partic-
ular solution is
x
p
= te
t
.
3.3.3 Damped resonance
view lecture
A more realistic study of resonance assumes an additional damping term. The
forced, damped harmonic oscillator equation may be written as
m x + x +kx = F cos t, (3.20)
where m > 0 is the oscillators mass, > 0 is the damping coecient, k >
0 is the spring constant, and F is the amplitude of the external force. The
homogeneous equation has characteristic equation
mr
2
+r +k = 0,
so that
r
=

2m

1
2m
_
2
4mk.
Since m, and k are all positive, the roots of the characteristic equation clearly
satisfy Re(r
) < 0. Therefore, both linearly independent homogeneous solutions

decay exponentially to zero, and the long-time asymptotic solution of (3.20)
reduces to the (non-decaying) particular solution. Since the initial conditions are
satised by the free constants multiplying the (decaying) homogeneous solutions,
the long-time asymptotic solution is independent of the initial conditions.
If we are only interested in the long-time asymptotic solution of (3.20), we
can proceed directly to the determination of a particular solution. As before,
we consider the complex ode
m z + z +kz = Fe
it
,
with x
p
= Re(z
p
). With the ansatz z
p
= Ae
it
, we have
m
2
A+iA+kA = F,
or
A =
F
(k m
2
) +i
.
To simplify, we dene
0
to be the natural frequency of the undamped oscillator,
so that
2
0
= k/m, and dene = /m, f = F/m. Therefore,
A =
f
(
2
0
2
) +i
=
_
f
(
2
0
2
)
2
+
2
_
_
(
2
0
2
) i
_
. (3.21)
To determine x
p
, we utilize the polar form of a complex number. The complex
number z = x +iy can be written in polar form as z = re
i
, where x = r cos ,
y = r sin, and r =
_
x
2
+y
2
, tan = y/x. We therefore write
(
2
0
2
) i = re
i
,
with
r =
_
(
2
0
2
)
2
+
2
, tan =

(
2
2
0
)
.
Using the polar form, A in (3.21) becomes
A =
_
f
_
(
2
0
2
)
2
+
2
_
e
i
,
and x
p
= Re(Ae
it
) becomes
x
p
=
_
f
_
(
2
0
2
)
2
+
2
_
Re
_
e
i(t+)
_
=
_
f
_
(
2
0
2
)
2
+
2
_
cos (t +). (3.22)
The particular solution given by (3.22), with
2
0
= k/m, = /m, f = F/m,
and tan = /m(
2
2
0
), is the long-time asymptotic solution of the forced,
damped, harmonic oscillator equation (3.20).
We conclude with a couple of observations about (3.22). First, if the forcing
frequency is equal to the natural frequency
0
of the undamped oscillator, then
A = iF/
0
, and x
p
= (F/
0
) sin
0
t. The oscillator position is observed
to be /2 out of phase with the external force, or in other words, the velocity
of the oscillator, not the position, is in phase with the force. Second, the value
of the forcing frequency
m
that maximizes the amplitude of oscillation is the
value of that minimizes the denominator of (3.22). To determine
m
we thus
minimize the function g(
2
) with respect to
2
, where
g(
2
) = (
2
0
2
)
2
+

2
2
m
2
.
Taking the derivative of g with respect to
2
and setting this to zero to determine
m
yields
2(
2
m
2
0
) +

2
m
2
= 0,
or
2
m
=
2
0

2
2m
2
.
We can interpret this result by saying that damping lowers the resonance
frequency of the undamped oscillator.
Chapter 4
The Laplace transform
The Laplace transform is most useful for solving linear, constant-coecient
odes when the inhomogeneous term or its derivative is discontinuous. Although
odes with discontinuous inhomogeneous terms can also be solved by adopting
already learned methods, we will see that the Laplace transform technique pro-
vides a simpler, more elegant solution.
4.1 Denitions and properties of the forward
and inverse Laplace transforms
view lecture
The main idea is to Laplace transform the constant-coecient dierential equa-
tion for x(t) into a simpler algebraic equation for the Laplace-transformed func-
tion X(s), solve this algebraic equation, and then transform X(s) back into
x(t). The correct denition of the Laplace transform and the properties that
this transform satises makes this solution method possible.
An exponential ansatz is used in solving homogeneous constant-coecient
odes, and the exponential function correspondingly plays a key role in dening
the Laplace transform. The Laplace transform of f(t), denoted by F(s) =
L{f(t)}, is dened by the integral transform
F(s) =
_

0
e
st
f(t)dt. (4.1)
The improper integral given by (4.1) diverges if f(t) grows faster than e
st
for
large t. Accordingly, some restriction on the range of s may be required to
guarantee convergence of (4.1), and we will assume without further elaboration
that these restrictions are always satised.
One of the useful properties of the Laplace transform is that it is a linear
operator. We have
L{c
1
f
1
(t) +c
2
f
2
(t)} =
_

0
e
st
(c
1
f
1
(t) +c
2
f
2
(t))dt
= c
1
_

0
e
st
f
1
(t)dt +c
2
_

0
e
st
f
2
(t)dt
= c
1
L{f
1
(t)} +c
2
L{f
2
(t)}.
45
46 CHAPTER 4. THE LAPLACE TRANSFORM
f(t) = L
1
{F(s)} F(s) = L{f(t)}
1. 1
1
s
2. e
at
1
s a
3. t
n
n!
s
n+1
4. sinat
a
s
2
+a
2
5. cos at
s
s
2
+a
2
6. sinhat
a
s
2
a
2
7. coshat
s
s
2
a
2
8. e
at
sinbt
b
(s a)
2
+b
2
9. e
at
cos bt
s a
(s a)
2
+b
2
10. t
n
e
at
n!
(s a)
n+1
11. u
c
(t)
e
cs
s
12. u
c
(t)f(t c) e
cs
F(s)
13. e
ct
f(t) F(s c)
14. (t c) e
cs
15. x(t) sX(s) x(0)
16. x(t) s
2
X(s) sx(0) x(0)
Table 4.1: Table of Laplace Transforms
4.1. DEFINITION AND PROPERTIES 47
Another useful property of the Laplace transformation is that there is a one-
to-one correspondence between functions and their Laplace transforms, so that
the inverse Laplace transform, also a linear operator, is well dened. This sug-
gests that construction of a table of Laplace transforms of commonly occurring
functions may be useful in nding both Laplace and inverse Laplace transforms.
Such a table is shown in Table 4.1 (and this table will be included with the nal
exam). In Table 4.1, n is a positive integer. Also, the cryptic entries for u
c
(t)
and (t c) will be explained later in 4.3.
The rows of Table 4.1 can be determined by a combination of direct in-
tegration and some tricks. We begin with a direct calculation of the Laplace
transform of 1 (line 1):
L{1} =
_

0
e
st
dt
=
1
s
e
st
|
0
=
1
s
.
We also compute directly the Laplace transform of e
ct
f(t) (line 13):
L{e
ct
f(t)} =
_

0
e
st
e
ct
f(t)dt
=
_

0
e
(sc)t
f(t)dt
= F(s c).
Now, the Laplace transform of e
at
(line 2) may be found using these two results:
L{e
at
} = L{e
at
1}
=
1
s a
.
The Laplace transform of sinat (line 4) may be found from the Laplace trans-
form of e
at
using a trick:
L{sinat} = Im{L{e
iat
}}
= Im{
1
s ia
}
= Im{
s +ia
s
2
+a
2
}
=
a
s
2
+a
2
.
Similarly, the Laplace transform of cos at (line 5) is
L{cos at} = Re{L{e
iat
}}
=
s
s
2
+a
2
.
The transforms of the hyperbolic functions can be found from the transform of
the exponential function. The transform of sinhat (line 6) is found using
L{sinhat} = {L{
1
2
(e
at
e
at
)}}
=
1
2
_
1
s a

1
s +a
_
=
a
s
2
a
2
;
and the transform of coshat (line 7) is found using
L{coshat} = {L{
1
2
(e
at
+e
at
)}}
=
1
2
_
1
s a
+
1
s +a
_
=
s
s
2
a
2
.
The transform of e
at
sinbt can be found using the transform of sinbt (line 4)
and line 13:
L{e
at
sinbt} =
b
(s a)
2
+b
2
,
and similarly for the transform of e
at
cos bt:
L{e
at
cos bt} =
s a
(s a)
2
+b
2
.
The transform of t
n
(line 3) requires another direct integration using succes-
sive integration by parts:
L{t
n
} =
_

0
e
st
t
n
dt
=
n
s
_

0
e
st
t
n1
dt
=
n(n 1)
s
2
_

0
e
st
t
n2
dt
= . . .
=
n!
s
n
_

0
e
st
dt
=
n!
s
n+1
.
The Laplace transform of t
n
e
at
can be found using the Laplace transform of t
n
together with line 13:
L{t
n
e
at
} =
n!
(s a)
n+1
.
We now transform the inhomogeneous constant-coecient, second-order, lin-
ear inhomogeneous ode for x = x(t),
a x +b x +cx = g(t),
4.2. SOLUTION OF INITIAL VALUE PROBLEMS 49
making use of the linearity of the Laplace transform:
aL{ x} +bL{ x} +cL{x} = L{g}.
To determine the Laplace transform of x(t) in terms of the Laplace transform
of x(t) and the initial conditions x(0) and x(0), we dene X(s) = L{x(t)}, and
integrate
_

0
e
st
xdt
by parts. We let
u = e
st
dv = xdt
du = se
st
dt v = x.
Therefore,
_

0
e
st
xdt = xe
st
|
0
+s
_

0
e
st
xdt
= sX(s) x(0),
where assumed convergence of the Laplace transform requires
lim
t
x(t)e
st
= 0.
Similarly, the Laplace transform of the second derivative is determined by inte-
grating
_

0
e
st
xdt
by parts and using the just derived result for the rst derivative. We let
u = e
st
dv = xdt
du = se
st
dt v = x,
so that
_

0
e
st
xdt = xe
st
|
0
+s
_

0
e
st
xdt
= x(0) +s (sX(s) x(0))
= s
2
X(s) sx(0) x(0),
where similarly we assume lim
t
x(t)e
st
= 0. We tabulate these important
results for the rst and second derivatives of x(t) in lines 15 and 16 of Table
4.1.
4.2 Solution of initial value problems
view lecture
We begin with a simple homogeneous ode and show that the Laplace transform
method yields an identical result to our previously learned method. We then
apply the Laplace transform method to solve an inhomogeneous equation.
Example: Solve x x2x = 0 with x(0) = 1 and x(0) = 0 by two dierent
methods.
The characteristic equation of the ode is determined from the ansatz x = e
rt
and is
r
2
r 2 = (r 2)(r + 1) = 0.
x(t) = c
1
e
t
+c
2
e
2t
.
To satisfy the initial conditions, we must have 1 = c
1
+ c
2
and 0 = c
1
+ 2c
2
,
requiring c
1
=
2
3
and c
2
=
1
3
. Therefore, the solution to the ode and initial
conditions is given by
x(t) =
2
3
e
t
+
1
3
e
2t
.
We now solve this example by the method of Laplace transform. Taking the
Laplace transform of both sides of the ode, using the linearity of the transform,
and applying our result for the transform of the rst and second derivatives, we
nd
[s
2
X(s) sx(0) x(0)] [sX(s) x(0)] [2X(s)] = 0,
or
X(s) =
(s 1)x(0) + x(0)
s
2
s 2
.
Note that the denominator of the right-hand-side is just the quadratic from the
characteristic equation of the homogeneous ode, and that this factor arises from
the derivatives of the exponential term in the Laplace transform integral.
Applying the initial conditions, we nd
X(s) =
s 1
(s 2)(s + 1)
. (4.2)
We have thus determined the Laplace transformed solution X(s) = L{x(t)} and
consequently need to compute the inverse Laplace transform x(t) = L
1
{X(s)}.
Now, direct inversion of (4.2) by searching Table 4.1 is not possible, but a
partial fraction expansion may be useful. In particular, we write
s 1
(s 2)(s + 1)
=
a
s 2
+
b
s + 1
,
from which we determine
s 1 = a(s + 1) +b(s 2);
or
a +b = 1, a 2b = 1,
with solution a =
1
3
and b =
2
3
. Therefore,
X(s) =
1
3

1
s 2
+
2
3

1
s + 1
,
and line 2 of Table 4.1 gives us the inverse transforms of each term separately
to yield
x(t) =
1
3
e
2t
+
2
3
e
t
,
identical to our previously determined solution.
4.3. HEAVISIDE AND DIRAC DELTA FUNCTIONS 51
Example: Solve x + x = sin2t with x(0) = 2 and x(0) = 1 by Laplace
transform methods.
Taking the Laplace transform of both sides of the ode, we nd
s
2
X(s) sx(0) x(0) +X(s) = L{sin2t}
=
2
s
2
+ 4
,
where the Laplace transform of sin2t made use of line 4 of Table 4.1. Substi-
tuting for x(0) and x(0) and solving for X(s), we obtain
X(s) =
2s
3
+s
2
+ 8s + 6
(s
2
+ 1)(s
2
+ 4)
.
To determine the inverse Laplace transform from Table 4.1, we apply a some-
what tedious partial fraction expansion:
2s
3
+s
2
+ 8s + 6
(s
2
+ 1)(s
2
+ 4)
=
as +b
s
2
+ 1
+
cs +d
s
2
+ 4
=
(as +b)(s
2
+ 4) + (cs +d)(s
2
+ 1)
(s
2
+ 1)(s
2
+ 4)
=
(a +c)s
3
+ (b +d)s
2
+ (4a +c)s + (4b +d)
(s
2
+ 1)(s
2
+ 4)
.
Equating coecients of powers of s in the numerators, we obtain two indepen-
dent systems of two equations and two unknowns:
a +c = 2, b +d = 1,
4a +c = 8, 4b +d = 6,
with solution a = 2, c = 0, b =
5
3
, and d =
2
3
, so that
X(s) =
2s
s
2
+ 1
+
5/3
s
2
+ 1

2/3
s
2
+ 4
.
From lines 4 and 5 of Table 4.1, we obtain the solution by taking inverse Laplace
transforms of the three terms separately, where a = 1 in the rst two terms,
and a = 2 in the third term:
x(t) = 2 cos t +
5
3
sint
1
3
sin2t.
4.3 Heaviside and Dirac delta functions
view lecture
view lecture
The Laplace transform technique becomes truly useful when solving odes with
discontinuous or impulsive inhomogeneous terms, these terms commonly mod-
eled using Heaviside or Dirac delta functions. We will discuss these functions
in turn, as well as their Laplace transforms.
4.3.1 Heaviside function
The Heaviside or unit step function, denoted here by u
c
(t), is dened by
u
c
(t) =
_
0, t < c;
1, t c.
For t < c, the value of u
c
(t) is zero and then steps up to one for t > c. The
reverse step function is obtained from 1u
c
(t), which is one for t < c and steps
down to zero for t c. The precise value of the Heaviside function at the single
point t = c is typically unimportant in applications. The Heaviside function can
also be used to construct a function which starts at zero, steps up to one over
some interval, say (a, b), and then steps back down to zero. A short reection
shows that the following combination of Heaviside functions has this property:
u
a
(t) u
b
(t) =
_
_
_
0, t < a;
1, a t < b;
0, t b.
(4.3)
The Laplace transform of the Heaviside function is determined by integra-
tion:
L{u
c
(t)} =
_

0
e
st
u
c
(t)dt
=
_

c
e
st
dt
=
e
cs
s
,
and is given in line 11 of Table 4.1.
The Heaviside function can be used to represent a translation of a function
f(t) a distance c in the positive direction. We have
u
c
(t)f(t c) =
_
0, t < c;
f(t-c), t c.
The Laplace transform is
L{u
c
(t)f(t c)} =
_

0
e
st
u
c
(t)f(t c)dt
=
_

c
e
st
f(t c)dt
=
_

0
e
s(t
+c)
f(t
)dt
= e
cs
_

0
e
st
f(t
)dt
= e
cs
L{f(t)},
where we have used the change of variable t
= t c. The translation of f(t) a

distance c in the positive t direction corresponds to the multiplication of F(s)
by the exponential e
cs
. This result is shown in line 12 of Table 4.1.
4.3. HEAVISIDE AND DIRAC DELTA FUNCTIONS 53
Piecewise discontinuous inhomogeneous terms can be modeled using Heav-
iside functions. For example, suppose the inhomogeneous term is represented
by a linearly increasing function, which then turns into a sinusoidal function for
t > t
. Explicitly,
f(t) =
_
at, if t < t
;
at
+b sin(t t
), if t t
.
We can rewrite f(t) as f(t) = at +g(t), with
g(t) =
_
0, if t < t
;
b sin(t t
) a(t t
), if t t
;
and using a Heaviside function, f(t) becomes
f(t) = at + [b sin(t t
) a(t t
)] u
t
(t)
= at +h(t t
)u
t
(t),
where we have dened
h(t) = b sint at.
Using line 12, the Laplace transform of f(t) is
F(s) = aL{t} +L{h(t t
)u
t
(t)}
= aL{t} +e
t
s
L{h(t)},
and where using lines 3 and 4,
L{t} =
1
s
2
,
and
L{h(t)} =
b
s
2
+
2

a
s
2
.
4.3.2 Dirac delta function
The Dirac delta function, denoted as (t), is dened by requiring that for any
function f(t),
_

f(t)(t)dt = f(0).
The usual view of the shifted Dirac delta function (t c) is that it is zero
everywhere except at t = c, at which point it is innite, and it integrates
to unity. The Dirac delta function is technically not a function, but is what
mathematicians call a distribution. Nevertheless, in most cases of practical
interest, it can be treated like a function, where physical results are obtained
following a nal integration.
There are many ways to represent the Dirac delta function as a limit of a
well-dened function. For our purposes, the most useful representation makes
use of the step-up, step-down function of (4.3):
(t c) = lim
0
1
2
(u
c
(t) u
c+
(t)).
Before taking the limit, the well-dened step-up, step-down function is zero
except over a small interval of width 2 centered at t = c, over which it takes
the large value 1/2. The integral of this function is one, independent of the
value of .
The Laplace transform of the Dirac delta function is easily found by inte-
gration using the denition of the delta function:
L{(t c)} =
_

0
e
st
(t c)dt
= e
cs
.
This result is shown in line 14 of Table 4.1.
4.4 Discontinuous or impulsive inhomogeneous
terms
view lecture
view lecture
We now solve some more challenging odes with discontinuous or impulsive
inhomogeneous terms.
Example: Solve 2 x + x + 2x = u
5
(t) u
20
(t), with x(0) = x(0) = 0
The inhomogeneous term here is a step-up, step-down function that is unity
over the interval (5, 20) and zero elsewhere. Taking the Laplace transform of
the ode using Table 4.1,
2(s
2
X(s) sx(0) x(0)) +sX(s) x(0) + 2X(s) =
e
5s
s

e
20s
s
.
Using the initial values and solving for X(s), we nd
X(s) =
e
5s
e
20s
s(2s
2
+s + 2)
.
To determine the solution for x(t) we now need to nd the inverse Laplace
transform. The exponential functions can be dealt with using line 12 of Table
4.1. We write
X(s) = (e
5s
e
20s
)H(s),
where
H(s) =
1
s(2s
2
+s + 2)
.
Then using line 12, we have
x(t) = u
5
(t)h(t 5) u
20
(t)h(t 20), (4.4)
where h(t) = L
1
{H(s)}.
To determine h(t) we need the partial fraction expansion of H(s). Since the
discriminant of 2s
2
+s + 2 is negative, we have
1
s(2s
2
+s + 2)
=
a
s
+
bs +c
2s
2
+s + 2
,
4.4. DISCONTINUOUS OR IMPULSIVE TERMS 55
yielding the equation
1 = a(2s
2
+s + 2) + (bs +c)s;
or after equating powers of s,
2a +b = 0, a +c = 0, 2a = 1,
yielding a =
1
2
, b = 1, and c =
1
2
. Therefore,
H(s) =
1/2
s

s +
1
2
2s
2
+s + 2
.
Inspecting Table 4.1, the rst term can be transformed using line 1, and the
second term can be transformed using lines 8 and 9, provided we complete the
square of the denominator and then massage the numerator. That is, rst we
complete the square:
2s
2
+s + 2 = 2
_
s
2
+
1
2
s + 1
_
= 2
_
(s +
1
4
)
2
+
15
16
_
;
and next we write
s +
1
2
2s
2
+s + 2
=
s +
1
2
2
_
(s +
1
4
)
2
+
15
16
_
=
_
1
2
_
_
s +
1
4
_
+
1
15
_
15
16
(s +
1
4
)
2
+
15
16
.
Therefore, the function H(s) can be written as
H(s) =
_
1
2
_
1
s

_
1
2
_
_
s +
1
4
_
(s +
1
4
)
2
+
15
16
_
1
2
15
_
_
15
16
(s +
1
4
)
2
+
15
16
.
The rst term is transformed using line 1, the second term using line 9 and the
third term using line 8. We nally obtain
h(t) =
1
2

1
2
e
t/4
cos (
15t/4)
_
1
2
15
_
e
t/4
sin(
15t/4), (4.5)
which when combined with (4.4) yields the rather complicated solution for x(t).
We briey comment that it is also possible to solve this example without
using the Laplace transform. Clearly from the form of the inhomogeneous term
and the initial conditions, x(t) = 0 for 0 t 5. We then solve the ode
between 5 t 20 with the inhomogeneous term equal to unity and initial
conditions x(5) = x(5) = 0. This requires nding the homogeneous solution,
particular solution, and solving for the two unknown constants. Finally, we solve
the homogeneous ode for t 20 using as boundary conditions the previously
determined values x(20) and x(20). The student may benet by trying this as
an exercise and attempting to obtain a nal solution that agrees with the form
given by (4.4) and (4.5).
Example: Solve 2 x + x + 2x = (t 5) with x(0) = x(0) = 0
Here the inhomogeneous term is an impulse at time t = 5. Taking the Laplace
transform of the ode using Table 4.1, and applying the initial conditions,
(2s
2
+s + 2)X(s) = e
5s
,
so that
X(s) =
e
5s
2s
2
+s + 2
=
1
2
e
5s
_
1
s
2
+
1
2
s + 1
_
=
1
2
e
5s
_
1
(s +
1
4
)
2
+
15
16
_
=
1
2
_
16
15
e
5s
_
_
_
15
16
(s +
1
4
)
2
+
15
16
_
_
.
The inverse Laplace transform may now be computed using lines 12 and 8 of
Table 4.1:
x(t) =
2
15
u
5
(t)e
(t5)/4
sin(
15(t 5)/4). (4.6)

It is of interest to solve this example without using a Laplace transform.
Clearly, x(t) = 0 up to the time of impulse at t = 5. Furthermore, after the
impulse the ode is homogeneous and can be solved with standard methods. The
only diculty is determining the initial conditions of the homogeneous ode at
t = 5. To do so, we integrate the ode across the singularity at t = 5:
2
_
5+
5
xdt +
_
5+
5
xdt + 2
_
5+
5
xdt =
_
5+
5
(t 5)dt.
Without proof, we assume the continuity of x(t) across t = 5. The integral over
the delta function is one and the derivatives of x can be integrated using the
fundamental theorem of calculus. Therefore,
2 ( x(5 +) x(5 )) + (x(5 +) x(5 )) + 4x(5) = 1;
and taking x(5) = 0, assuming x(t) is continuous at t = 5, and x(5 ) = 0, we
obtain as 0:
x(5) =
1
2
,
which is the appropriate initial condition on the derivative, as can be seen by
taking the derivative of our solution (4.6), assuming t > 5, and then taking the
limit as t 5 from above. Notice that x is discontinuous at t = 5, jumping
from zero to one-half as t crosses 5.
Chapter 5
Series solutions of
second-order linear
homogeneous dierential
equations
view lecture
We consider the second-order linear homogeneous dierential equation for y =
y(x):
P(x)y
+Q(x)y
+R(x)y = 0, (5.1)
where P(x), Q(x) and R(x) are polynomials in x, with no common polynomial
factors (that could be divided out). The value x = x
0
is called an ordinary
point of (5.1) if P(x
0
) = 0, and is called a singular point if P(x
0
) = 0. Singular
points will later be further classied as regular singular points and irregular
singular points. Our goal is to nd two independent solutions of (5.1), valid in
a neighborhood about x = x
0
.
5.1 Ordinary points
If x
0
is an ordinary point of (5.1), then it is possible to determine two power
series solutions, or Taylor series, for y = y(x) that converge in a neighborhood
of x = x
0
. We illustrate the method of solution by solving two examples.
Example: Find the general solution of y
+y = 0.
By now, you should know that the general solution is y(x) = a
0
cos x +a
1
sinx,
with a
0
and a
1
constants. To nd a power series solution about the point x = 0,
we write
y(x) =
n=0
a
n
x
n
;
57
58 CHAPTER 5. SERIES SOLUTIONS
and upon dierentiating term-by-term
y
(x) =
n=1
na
n
x
n1
,
and
y
(x) =
n=2
n(n 1)a
n
x
n2
.
Substituting the power series for y and its derivatives into the dierential equa-
tion to be solved, we obtain
n=2
n(n 1)a
n
x
n2
+
n=0
a
n
x
n
= 0. (5.2)
The power-series solution method requires combining the two sums on the left-
hand-side of (5.2) into a single power series in x. To shift the exponent of x
n2
in the rst sum upward by two to obtain x
n
, we need to shift the summation
index downward by two, i.e.,
n=2
n(n 1)a
n
x
n2
=
n=0
(n + 2)(n + 1)a
n+2
x
n
.
We can then combine the two sums in (5.2) to obtain
n=0
{(n + 2)(n + 1)a
n+2
+a
n
} x
n
= 0. (5.3)
For (5.3) to be satised, the coecient of each power of x must vanish separately.
(This can be proved by setting x = 0 after successive dierentiation.) We
therefore obtain the recurrence relation
a
n+2
=
a
n
(n + 2)(n + 1)
.
We note that even and odd coecients decouple. We thus obtain two indepen-
dent sequences starting with rst term a
0
or a
1
. Developing these sequences,
we have for the sequence beginning with a
0
:
a
0
,
a
2
=
1
2
a
0
,
a
4
=
1
4 3
a
2
=
1
4 3 2
a
0
,
a
6
=
1
6 5
a
4
=
1
6!
a
0
;
and the general coecient in this sequence for n = 0, 1, 2, . . . is
a
2n
=
(1)
n
(2n)!
a
0
.
5.1. ORDINARY POINTS 59
Also, the sequence beginning with a
1
is
a
1
,
a
3
=
1
3 2
a
1
,
a
5
=
1
5 4
a
3
=
1
5 4 3 2
a
1
,
a
7
=
1
7 6
a
5
=
1
7!
a
1
;
and the general coecient in this sequence for n = 0, 1, 2, . . . is
a
2n+1
=
(1)
n
(2n + 1)!
a
1
.
Using the principle of superposition, the general solution is therefore
y(x) = a
0
n=0
(1)
n
(2n)!
x
2n
+a
1
n=0
(1)
n
(2n + 1)!
x
2n+1
= a
0
_
1
x
2
2!
+
x
4
4!
. . .
_
+a
1
_
x
x
3
3!
+
x
5
5!
. . .
_
= a
0
cos x +a
1
sinx,
as it should be.
We now solve the Airys Equation in our next example.
Example: Find the general solution of y
xy = 0.
view lecture
With
y(x) =
n=0
a
n
x
n
,
the dierential equation becomes
n=2
n(n 1)a
n
x
n2
n=0
a
n
x
n+1
= 0. (5.4)
We shift both sums to x
n
by shifting the exponent in the rst sum up by two
and in the second sum down by one, i.e.,
n=2
n(n 1)a
n
x
n2
=
n=0
(n + 2)(n + 1)a
n+2
x
n
;
n=0
a
n
x
n+1
=
n=1
a
n1
x
n
.
When combining the two sums, we note the extra n = 0 term in the rst sum
and write this term, 2a
2
x
0
, explicitly. Therefore, (5.4) becomes
2a
2
+
n=1
{(n + 2)(n + 1)a
n+2
a
n1
} x
n
= 0. (5.5)
Setting coecients of powers of x to zero, we obtain a
2
= 0, and the recursion
relation
a
n+2
=
1
(n + 2)(n + 1)
a
n1
. (5.6)
Three sequences of coecients, starting with a
0
, a
1
and a
2
decouple. In partic-
ular the sequences are
a
0
, a
3
, a
6
, a
9
, . . . ;
a
1
, a
4
, a
7
, a
10
, . . . ;
a
2
, a
5
, a
8
, a
11
. . . .
Since a
2
= 0, we have for the last sequence
a
2
= a
5
= a
8
= a
11
= . . . = 0.
We compute the rst four nonzero terms in the power series with coecients
corresponding to the rst two sequences. Starting with a
0
, we have
a
0
,
a
3
=
1
3 2
a
0
,
a
6
=
1
6 5 3 2
a
0
,
a
9
=
1
9 8 6 5 3 2
a
0
;
and starting with a
1
,
a
1
,
a
4
=
1
4 3
a
1
,
a
7
=
1
7 6 4 3
a
1
,
a
10
=
1
10 9 7 6 4 3
a
1
.
The general solution for y = y(x), can therefore be written as
y(x) = a
0
_
1 +
x
3
6
+
x
6
180
+
x
9
12960
+. . .
_
+a
1
_
x +
x
4
12
+
x
7
504
+
x
10
45360
+. . .
_
= a
0
y
0
(x) +a
1
y
1
(x).
Suppose we would like to graph the solutions y = y
0
(x) and y = y
1
(x)
versus x by solving the dierential equation y
xy = 0 numerically. What
initial conditions should we use? Clearly, y = y
0
(x) solves the ode with initial
values y(0) = 1 and y
(0) = 0, while y = y
1
(x) solves the ode with initial values
y(0) = 0 and y
(0) = 1. The numerical solutions, obtained using MATLAB,

are shown in Fig. 5.1. Note that the solutions oscillate for negative x and
grow exponentially for positive x. This can be understood by recalling that
y
+y = 0 has oscillatory sine and cosine solutions and y
y = 0 has exponential
hyperbolic sine and cosine solutions.
5.2. REGULAR SINGULAR POINTS: EULER EQUATIONS 61
10 8 6 4 2 0 2
2
1
0
1
2
y
0
Airys functions
10 8 6 4 2 0 2
2
1
0
1
2
x
y
1
Figure 5.1: Numerical solution of Airys equation.
5.2 Regular singular points: Euler equations
view lecture
view lecture
The value x = x
0
is called a regular singular point of the ode
(x x
0
)
2
y
+p(x)(x x
0
)y
+q(x)y = 0, (5.7)
if p(x) and q(x) have convergent Taylor series about x = x
0
, i.e., p(x) and q(x)
can be written as a power-series in (x x
0
):
p(x) = p
0
+p
1
(x x
0
) +p
2
(x x
0
)
2
+. . . ,
q(x) = q
0
+q
1
(x x
0
) +q
2
(x x
0
)
2
+. . . ,
with p
0
, p
1
, p
2
, q
0
, q
1
, q
2
, etc. constants (all or some of which may be zero). Any
point x = x
0
that is not an ordinary point or a regular singular point is called an
irregular singular point. Many important dierental equations of physical inter-
est associated with special functions named after now famous mathematicians
like Bessel, Legendre, Hermite, Laguerre and Chebyshev have regular singular
points.
Here, we will only consider the simplest ode with a regular singular point at
x = 0. This ode is called an Euler equation, and has the form
x
2
y
+xy
+y = 0, (5.8)
with and constants. Note that (5.7) reduces to an Euler equation (about
x = x
0
) when one considers only the leading-order term in the Taylor series
expansion of the functions p(x) and q(x). In fact, taking p(x) = p
0
and q(x) = q
0
and solving the associated Euler equation results in the leading-order solution
to the more general ode (5.7). Often, this is sucient to obtain initial conditions
for numerical solution of the full ode. Students wishing to learn how to nd the
general solution of (5.7) can consult the textbook.
An appropriate ansatz for (5.8) is y = x
r
, when x > 0 and y = (x)
r
when
x < 0, (or more generally, y = |x|
r
for all x), with r constant. After substitution
into (5.8), we obtain for both positive and negative x
r(r 1)|x|
r
+r|x|
r
+|x|
r
= 0,
and we observe that our ansatz is rewarded by cancellation of |x|
r
. We thus
obtain the following quadratic equation for r:
r
2
+ ( 1)r + = 0, (5.9)
which can be solved using the quadratic formula. Three cases immediately
appear: (i) real distinct roots, (ii) complex conjugate roots, (iii) repeated roots.
Students may recall being in a similar situation when solving the second-order
linear homogeneous ode with constant coecients. Indeed, it is possible to
directly transform the Euler equation into an equation with constant coecients
so that our previous results can be used.
The appropriate transformation for x > 0 is to let x = e
and y(x) = Y (),

so that the ansatz y(x) = x
r
becomes the ansatz Y () = e
r
, appropriate if Y ()
satises a constant coecient ode. If x < 0, then the appropriate transformation
is x = e
, since e
> 0. We need only consider x > 0 here and subsequently

generalize our result by replacing x everywhere by its absolute value.
We thus transform the dierential equation (5.8) for y = y(x) into a dif-
ferential equation for Y = Y (), using x = e
and = lnx. By the chain

rule,
dy
dx
=
dY
d
d
dx
=
1
x
dY
d
= e
dY
d
,
so that
d
dx
= e
d
d
.
The second derivative becomes
d
2
y
dx
2
= e
d
d
_
e
dY
d
_
= e
2
_
d
2
Y
d
2

dY
d
_
.
Upon substitution of the derivatives of y into (5.8), and using x = e
, we obtain
e
2
_
e
2
(Y
)
_
+e
_
e
_
+Y = Y
+ ( 1)Y
+Y
= 0.
5.2. REGULAR SINGULAR POINTS: EULER EQUATIONS 63
As expected, the ode for Y = Y () has constant coecients, and with Y = e
r
,
the characteristic equation for r is given by (5.9). We now directly transfer
previous results obtained for the constant coecient second-order linear homo-
geneous ode.
5.2.1 Real, distinct roots
This simplest case needs no transformation. If ( 1)
2
4 > 0, then with r
the real roots of (5.9), the general solution is

y(x) = c
1
|x|
r
+
+c
2
|x|
r
.
5.2.2 Complex conjugate roots
If ( 1)
2
4 < 0, we can write the complex roots of (5.9) as r
= i.
Recall the general solution for Y = Y () is given by
Y () = e
(Acos +Bsin) ;
and upon transformation, and replacing x by |x|,
y(x) = |x|
(Acos (ln|x|) +Bsin(ln|x|)) .

5.2.3 Repeated roots
If ( 1)
2
4 = 0, there is one real root r of (5.9). The general solution for
Y is
Y () = (c
1
+c
2
) e
r
,
yielding
y(x) = (c
1
+c
2
ln|x|) |x|
r
.
We now give examples illustrating these three cases.
Example: Solve 2x
2
y
+ 3xy
y = 0 for 0 x 1 with two-point

boundary condition y(0) = 0 and y(1) = 1.
Since x > 0, we try y = x
r
and obtain the characteristic equation
2r(r 1) + 3r 1 = 0,
or
2r
2
+r 1 = (2r 1)(r + 1) = 0.
Since the characteristic equation has two real roots, the general solution is given
by
y(x) = c
1
x
1
2
+c
2
x
1
.
We now encounter for the rst time two-point boundary conditions, which can
be used to determine the coecients c
1
and c
2
. Since y(0)=0, we must have
c
2
= 0. Applying the remaining condition y(1) = 1, we obtain the unique
solution
y(x) =
x.
Note that x = 0 is called a singular point of the ode since the general solution
is singular at x = 0 when c
2
= 0. Our boundary condition imposes that y(x) is
nite at x = 0 removing the singular solution. Nevertheless, y
remains singular
at x = 0. Indeed, this is why we imposed a two-point boundary condition rather
than specifying the value of y
(0) (which is innite).

Example: Find the general solution of x
2
y
+xy
+y = 0 for x > 0.
With the ansatz y = x
r
, we obtain
0 = r(r 1) +r + 1
= r
2
+ 1,
so that r = i. Therefore, with = lnx, we have Y () = Acos +Bsin, and
our solution for y(x) is
y(x) = Acos (lnx) +Bsin(lnx).
2
y
+ 5xy
+ 4y = 0 for x > 0.
With the ansatz y = x
r
, we obtain
0 = r(r 1) + 5r + 4
= r
2
+ 4r + 4
= (r + 2)
2
,
so that there is a repeated root r = 2. With = lnx, we have Y () =
(c
1
+c
2
)e
2
, so that
y(x) =
c
1
+c
2
lnx
x
2
.
Chapter 6
Systems of rst-order linear
equations
6.1 Determinants and the eigenvalue problem
view lecture
We begin by reviewing some basic linear algebra. For the simplest 2 2 case,
let
A =
_
a b
c d
_
, x =
_
x
1
x
2
_
,
and consider the homogeneous equation
Ax = 0. (6.1)
When does there exist a nontrivial (not identically zero) solution for x? To
answer this question, we solve directly the system
ax
1
+bx
2
= 0,
cx
1
+dx
2
= 0.
Multiplying the rst equation by d and the second by b, and subtracting the
second equation from the rst, results in
(ad bc)x
1
= 0.
Similarly, multiplying the rst equation by c and the second a, and subtracting
the rst equation from the second, results in
(ad bc)x
2
= 0.
Therefore, a nontrivial solution of (6.1) exists only if ad bc = 0. If we dene
the determinant of the 22 matrix A to be det (A) = adbc, then we say that
a nontrivial solution to (6.1) exists provided det (A) = 0.
The same calculation may be repeated for a 3 3 matrix. If
A =
_
_
a b c
d e f
g h i
_
_
, x =
_
_
x
1
x
2
x
3
_
_
,
65
66 CHAPTER 6. SYSTEMS OF EQUATIONS
then there exists a nontrivial solution to (6.1) provided det (A) = 0, where
det (A) = a(ei fh)b(di fg)+c(dheg). The denition of the determinant
can be further generalized to any nn matrix, and is typically studied in a rst
course on linear algebra.
We now consider the eigenvalue problem. For A an n n matrix and x an
n 1 column vector, the eigenvalue problem solves the equation
Ax = x (6.2)
for eigenvalues
i
and corresponding eigenvectors x
i
. We rewrite the eigenvalue
equation (6.2) as
(AI)x = 0, (6.3)
where I is the nn identity matrix. A nontrivial solution of (6.3) exists provided
det (AI) = 0. (6.4)
Equation (6.4) is an nth order polynomial equation in , and is called the
characteristic equation of A. The characteristic equation can be solved for
the eigenvalues, and for each eigenvalue, a corresponding eigenvector can be
determined directly from (6.2).
6.2 Two coupled rst-order linear homogeneous
dierential equations
view lecture
view lecture
With A a 2 2 constant matrix and x a 2 1 column vector, we consider the
system of dierential equations given by
x = Ax. (6.5)
We will consider three cases separately by example: (i) eigenvalues of A are real
and there are two linearly independent eigenvectors; (ii) eigenvalues of A are
complex conjugates, and; (iii) A has only one linearly independent eigenvector.
The three cases to be considered may remind students of our earlier solution
to the second-order linear (constant-coecient) homogeneous equation, and in-
deed, a system of two rst-order linear homogeneous equations can be converted
into a single second-order linear homogeneous equation and solved by means of a
characteristic equation. We will illustrate this equivalence. Nevertheless, view-
ing the problem as a system of rst-order equations introduces the important
concept of the phase space, and also has some interesting physical applications
such as normal modes of oscillation and linear stability.
6.2.1 Two distinct real eigenvalues
We illustrate the solution method by example.
6.2. COUPLED FIRST-ORDER EQUATIONS 67
1
= x
1
+x
2
, x
2
= 4x
1
+x
2
.
The equation to be solved may be rewritten in matrix form as
d
dt
_
x
1
x
2
_
=
_
1 1
4 1
__
x
1
x
2
_
,
or using vector notation, written as (6.5). We take as our ansatz x = ve
t
,
where v and are independent of t. Upon substitution into (6.5), we obtain
ve
t
= Ave
t
;
and upon cancellation of the exponential, we obtain the eigenvalue problem
Av = v. (6.6)
Solving the eigenvalue problem, we have
0 = det (AI)
=
1 1
4 1
= (1 )
2
4
=
2
2 3
= ( 3)( + 1).
Therefore, the two eigenvalues are
1
= 3 and
2
= 1. To determine the
corresponding eigenvectors, we substitute the eigenvalues successively into (A
I)v = 0. For
1
= 3, and eigenvector v
1
= (v
11
, v
21
)
T
, we have
2v
11
+v
21
= 0,
4v
11
2v
21
= 0.
Clearly, the second equation is just the rst equation multiplied by 2, so only
one equation is linearly independent. This will always be true, so for the 2 2
case we need only consider the rst row of the matrix. The rst eigenvector
therefore satises v
21
= 2v
11
. Recall that an eigenvector is only unique up
to multiplication by a constant. We may take v
11
= 1 for convenience. For
2
= 1, and eigenvector v
2
= (v
12
, v
22
)
T
, we have
2v
12
+v
22
= 0,
so that v
22
= 2v
12
. Here, we take v
12
= 1. Therefore, our eigenvalues and
eigenvectors are given by
1
= 3, v
1
=
_
1
2
_
;
2
= 1, v
2
=
_
1
2
_
.
Using the principle of linear superposition, the general solution to the ode is
therefore
x = c
1
v
1
e
1
t
+c
2
v
2
e
2
t
,
or more specically,
x
1
= c
1
e
3t
+c
2
e
t
,
x
2
= 2c
1
e
3t
2c
2
e
t
.
2 1.5 1 0.5 0 0.5 1 1.5 2
2
1.5
1
0.5
0
0.5
1
1.5
2
x
1
x
2
phase space diagram
Figure 6.1: Phase space diagram for example with two real eigenvalues of oppo-
site sign.
It is useful to visualize the solution in the phase space, with x-axis x
1
and
y-axis x
2
. We show the phase space plot in Fig. 6.1. The dark lines represent
the direction of the eigenvectors. If c
2
= 0, the motion is along the eigenvector
v
1
with x
2
= 2x
1
and motion is away from the origin (arrows pointing out) since
the eigenvalue
1
= 3 > 0. If c
1
= 0, the motion is along the eigenvector v
2
with x
2
= 2x
1
and motion is towards the origin (arrows pointing in) since the
eigenvalue
2
= 1 < 0. When the eigenvalues are real and of opposite signs,
the origin is called a saddle point. Almost all trajectories (with the exception
of those with initial conditions exactly satisfying x
2
(0) = 2x
1
(0)) eventually
move away from the origin as t increases.
The current example can also be solved by converting the system of two
rst-order equations into a single second-order equation. Considering again the
system of equations
x
1
= x
1
+x
2
,
x
2
= 4x
1
+x
2
,
we dierentiate the rst equation and proceed to eliminate x
2
as follows:
x
1
= x
1
+ x
2
= x
1
+ 4x
1
+x
2
= x
1
+ 4x
1
+ x
1
x
1
= 2 x
1
+ 3x
1
.
Therefore, the equivalent second-order linear homogeneous equation is given by
x
1
2 x
1
3x
1
= 0,
which may be solved by the usual method. Initial conditions x
1
(0) and x
2
(0)
convert into initial conditions x
1
(0) and x
1
(0) = x
1
(0) + x
2
(0). In general, a
system of n rst-order linear homogeneous equations can be converted into an
equivalent n-th order linear homogeneous equation. Numerical methods usually
require the conversion in reverse, that is, a conversion of an n-th order equation
into a system of n rst-order equations.
1
= 3x
1
+
2x
2
, x
2
=
2x
1
2x
2
.
The equations in matrix form are
d
dt
_
x
1
x
2
_
=
_
3

2
2 2
__
x
1
x
2
_
.
The ansatz x = ve
t
leads to the eigenvalue problem Av = v, with A the
matrix above. The eigenvalues are determined from
0 = det (AI)
=
2 2
= (3 +)(2 +) 2
=
2
+ 5 + 4
= ( + 4)( + 1).
Therefore, the eigenvalues of A are
1
= 4,
2
= 1. Proceeding to determine
the associated eigenvectors, for
1
= 4:
v
11
+
2v
21
= 0,
and for
2
= 1:
2v
12
+
2v
22
= 0.
Taking the normalization v
11
= 1 and v
12
= 1, we obtain for the eigenvalues
and associated eigenvectors
1
= 4, v
1
=
_
1
2/2
_
;
2
= 1, v
2
=
_
1
2
_
.
The general solution to the ode is therefore
_
x
1
x
2
_
= c
1
_
1
2/2
_
e
4t
+c
2
_
1
2
_
e
t
.
We show the phase space plot in Fig. 6.2. If c
2
= 0, the motion is along
the eigenvector v
1
with x
2
= (
2/2)x
1
with eigenvalue
1
= 4 < 0. If
c
1
= 0, the motion is along the eigenvector v
2
with x
2
=
2x
1
with eigenvalue
2
= 1 < 0. When the eigenvalues are real and of same signs, the origin is
called a node. A node may be attracting or repelling depending on whether the
eigenvalues are both negative (as is the case here) or positive. Observe that the
trajectories collapse onto the v
2
eigenvector since
1
<
2
< 0 and decay is
more rapid along the v
1
direction.
2 1.5 1 0.5 0 0.5 1 1.5 2
2
1.5
1
0.5
0
0.5
1
1.5
2
x
1
x
2
phase space diagram
Figure 6.2: Phase space diagram for example with two real eigenvalues of same
sign.
6.2.2 Complex conjugate eigenvalues
1
=
1
2
x
1
+x
2
, x
2
= x
1
1
2
x
2
.
d
dt
_
x
1
x
2
_
=
_

1
2
1
1
1
2
__
x
1
x
2
_
.
The ansatz x = ve
t
leads to the equation
0 = det (AI)
=
1
2
1
1
1
2

= ( +
1
2
)
2
+ 1.
Therefore, (+
1
2
)
2
= 1, or
+
1
2
= i, or
=
1
2
i; and we observe that
the eigenvalues occur as a complex conjugate pair. We denote
+
= =
1
2
+i
and
=

=
1
2
i. The eigenvectors also occur as a complex conjugate
pair, since for A a real matrix, if Av = v, then A v =

v. The eigenvector
v associated with eigenvalue satises iv
1
+ v
2
= 0, and normalizing with
v
1
= 1, we have
=
1
2
+i, v =
_
1
i
_
.
The general solution to the ode is therefore
x = c
1
ve
t
+c
2
ve
t
;
2 1.5 1 0.5 0 0.5 1 1.5 2
2
1.5
1
0.5
0
0.5
1
1.5
2
x
1
x
2
phase space diagram
Figure 6.3: Phase space diagram for example with complex conjugate eigenval-
ues.
and since x is real so that x = x, we require c
2
= c
1
. Therefore, we may write
x = Re
_
cve
t
_
,
where c is in general a complex constant. Making use of our eigenvalue and
eigenvector, and writing for aesthetics c = AiB, we have
x = Re
_
c
_
1
i
_
e
(
1
2
+i)t
_
= e
1
2
t
Re
_
(AiB)
_
1
i
_
(cos t +i sint)
_
= e
1
2
t
_
Acos t +Bsint
Asint +Bcos t
_
= Ae
1
2
t
_
cos t
sint
_
+Be
1
2
t
_
sint
cos t
_
.
The corresponding phase space diagram is shown in Fig. 6.3. We say the origin
is a spiral point. If the real part of the complex eigenvalue is negative, as is
the case here, then all solutions spiral into the origin. If the real part of the
eigenvalue is positive, then solutions spiral out of the origin.
The direction of the spiralhere, it is clockwisecan be determined using
techniques from vector calculus. If a particle of unit mass moves along a phase
space trajectory, then the angular momentum of the particle about the origin
is equal to the cross product of the position and velocity vectors: L = x x.
With both the position and velocity vectors lying in the two-dimensional phase
space plane, the angular momentum vector is perpendicular to this plane. With
x = (x
1
, x
2
, 0), x = ( x
1
, x
2
, 0),
then
L = (0, 0, L), with L = x
1
x
2
x
2
x
1
.
By the right-hand-rule, a clockwise rotation corresponds to L < 0 and a coun-
terclockwise rotation to L > 0.
Computation of L for our current example proceeds from the dierential
equations:
x
1
x
2
x
2
x
1
= x
1
(x
1
1
2
x
2
) x
2
(
1
2
x
1
+x
2
)
= (x
2
1
+x
2
2
)
< 0;
and since L < 0, the spiral rotation is clockwise.
6.2.3 Repeated eigenvalues with one eigenvector
view lecture
1
= x
1
x
2
, x
2
= x
1
+ 3x
2
.
d
dt
_
x
1
x
2
_
=
_
1 1
1 3
__
x
1
x
2
_
. (6.7)
The ansatz x = ve
t
leads to the equation
0 = det (AI)
=
1 1
1 3
= (1 )(3 ) + 1
=
2
4 + 4
= ( 2)
2
.
Therefore, = 2 is a repeated eigenvalue. The associated eigenvector is found
from v
1
v
2
= 0, or v
2
= v
1
; and normalizing with v
1
= 1, we have
= 2, v =
_
1
1
_
.
We have thus found a single solution to the ode, given by
x(t) = c
1
_
1
1
_
e
2t
,
and we need to determine a second linearly independent solution in order to be
able to satisfy the initial conditions. An ansatz of t times the rst solution is
tempting, but in fact will not succeed. Here, we determine a second solution by
converting the two rst-order equations to a second-order equation, determine
the solution of the second-order equation using previously learned methods, and
then convert back to a solution for x
1
and x
2
. Accordingly, starting with the
rst-order equations
x
1
= x
1
x
2
,
x
2
= x
1
+ 3x
2
,
we have
x
1
= x
1
x
2
= x
1
(x
1
+ 3x
2
)
= x
1
x
1
3(x
1
x
1
)
= 4 x
1
4x
1
.
Therefore, the second-order ode for x
1
is given by
x
1
4 x
1
+ 4x
1
= 0.
With the ansatz x = e
rt
, the characteristic equation is r
2
4r + 4 = 0, or
(r 2)
2
= 0, so as one would expect, there is a repeated root r = 2. The general
solution for x
1
is therefore
x
1
= (c
1
+tc
2
)e
2t
.
To determine x
2
, we rst nd
x
1
= (2c
1
+ (1 + 2t)c
2
) e
2t
,
so that
x
2
= x
1
x
1
= (c
1
+tc
2
)e
2t
(2c
1
+ (1 + 2t)c
2
) e
2t
= c
1
e
2t
+c
2
(1 t)e
2t
.
Therefore, combining our results for x
1
and x
2
, we have
_
x
1
x
2
_
= c
1
_
1
1
_
e
2t
+c
2
__
0
1
_
+
_
1
1
_
t
_
e
2t
.
Our missing linearly independent solution is thus determined to be
x(t) = c
2
__
0
1
_
+
_
1
1
_
t
_
e
2t
. (6.8)
The second term of (6.8) is seen to be t times the rst solution; however, this
is not sucient. Indeed, the correct ansatz to nd the second solution directly
is
x = (w +tv) e
t
, (6.9)
where and v is the eigenvalue and eigenvector of the rst solution, and w
is an unknown vector to be determined. To illustrate this direct method, we
substitute (6.9) into x = Ax, assuming Av = v . Cancelling the exponential,
we obtain
v +(w +tv) = Aw +tv.
2 1.5 1 0.5 0 0.5 1 1.5 2
2
1.5
1
0.5
0
0.5
1
1.5
2
x
1
x
2
phase space diagram
Figure 6.4: Phase space diagram for example with only one eigenvector.
Further canceling the common term tv and rewriting yields
(AI) w = v. (6.10)
If A has only a single linearly independent eigenvector v, then (6.10) can be
solved for w (otherwise, it can not). Using A, and v of our present example,
(6.10) is the system of equations given by
_
1 1
1 1
__
w
1
w
2
_
=
_
1
1
_
.
The rst and second equation are the same, so that w
2
= (w
1
+1). Therefore,
w =
_
w
1
(w
1
+ 1)
_
= w
1
_
1
1
_
+
_
0
1
_
.
Notice the rst term repeats the rst found solution, i.e., a constant times the
eigenvector, and the second term is new. We therefore take w
1
= 0 and obtain
w =
_
0
1
_
,
as before.
The phase space diagram for this ode is shown in Fig. 6.4. The dark line is
the single eigenvector v of the matrix A. When there is only a single eigenvector,
the origin is called an improper node.
There is a denite counterclockwise rotation to the phase space trajectories,
and this can be conrmed from the calculation
x
1
x
2
x
2
x
1
= x
1
(x
1
+ 3x
2
) x
2
(x
1
x
2
)
6.3. NORMAL MODES 75
= x
2
1
+ 2x
1
x
2
+x
2
2
= (x
1
+x
2
)
2
> 0.
6.3 Normal modes
view lecture
We now consider an application of the eigenvector analysis to the coupled mass-
spring system shown in Fig. 6.5. The position variables x
1
and x
2
are measured
from the equilibrium positions of the masses. Hookes law states that the spring
force is linearly proportional to the extension length of the spring, measured
from equilibrium. By considering the extension of the spring and the sign of the
force, we write Newtons law F = ma separately for each mass:
m x
1
= kx
1
k
12
(x
1
x
2
),
m x
2
= kx
2
k
12
(x
2
x
1
).
Further rewriting by collecting terms proportional to x
1
and x
2
yields
m x
1
= (k +k
12
)x
1
+k
12
x
2
,
m x
2
= k
12
x
1
(k +k
12
)x
2
.
The equations for the coupled mass-spring system form a system of two second-
order linear homogeneous odes. In matrix form, m x = Ax, or explicitly,
m
d
2
dt
2
_
x
1
x
2
_
=
_
(k +k
12
) k
12
k
12
(k +k
12
)
__
x
1
x
2
_
. (6.11)
In analogy to a system of rst-order equations, we try the ansatz x = ve
t
, and
upon substitution into (6.11) we obtain the eigenvalue problem Av = m
2
v.
The values of m
2
are determined by solving the characteristic equation
0 = det (Am
2
I)
=
(k +k
12
) m
2
k
12
k
12
(k +k
12
) m
2
= (m
2
+k +k
12
)
2
k
2
12
.
The equation for m
2
is thus
m
2
= k k
12
k
12
,
and the solution for the two eigenvalues is therefore
2
1
= k/m,
2
2
= (k + 2k
12
)/m.
Since
2
1
,
2
2
< 0, both values of are pure imaginary, and x
1
(t) and x
2
(t)
oscillate with angular frequencies
1
= |
1
| and
2
= |
2
|, where
1
=
_
k/m,
2
=
_
(k + 2k
12
)/m.
Figure 6.5: Coupled harmonic oscillators (gure from Eric Weissteins World
of Physics).
The positions of the oscillating masses in general contain time dependencies of
the form sin
1
t, cos
1
t, and sin
2
t, cos
2
t.
It is of further interest to determine the eigenvectors, or so-called normal
modes of oscillation, associated with the two distinct angular frequencies. With
specic initial conditions corresponding to an eigenvector, the mass will oscillate
with a single frequency. The eigenvector associated with m
2
1
satises
k
12
v
11
+k
12
v
12
= 0,
so that v
11
= v
12
. The normal mode associated with the frequency
1
=
_
k/m thus follows a motion where x
1
= x
2
. Referring to Fig. 6.5, during this
motion the center spring length does not change, and the two masses oscillate
as if the center spring were absent (which is why the frequency of oscillation is
independent of k
12
).
Next, we determine the eigenvector associated with m
2
2
:
k
12
v
21
+k
12
v
22
= 0,
so that v
21
= v
22
. The normal mode associated with the frequency
2
=
_
(k + 2k
12
)/m thus follows a motion where x
1
= x
2
. Again referring to
Fig. 6.5, during this motion the two masses simultaneously press on or pull
against each side of the middle spring.
A general solution for x(t) can be constructed from the eigenvalues and
eigenvectors. Our ansatz was x = ve
t
, and we can linearly superimpose our
solutions with the additional requirement that x is real. The nal result, with
1
=
_
k/m and
2
=
_
(k + 2k
12
)/m, is
_
x
1
x
2
_
=
_
1
1
_
(Acos
1
t +Bsin
1
t) +
_
1
1
_
(C cos
2
t +Dsin
2
t) ,
where the now real constants A, B, C, and D can be determined from the four
independent initial conditions, x
1
(0), x
2
(0), x
1
(0), and x
2
(0).
Chapter 7
Nonlinear dierential
equations and bifurcation
theory
Reference: Strogatz, Sections 2.2, 2.4, 3.1, 3.2, 3.4, 6.3, 6.4, 8.2
view lecture
We now turn our attention to nonlinear dierential equations. In particular, we
study how small changes in the parameters of a system can result in qualita-
tive changes in the dynamics. These qualitative changes in the dynamics are
called bifurcations. To understand bifurcations, we rst need to understand the
concepts of xed points and stability.
7.1 Fixed points and stability
Consider the one dimensional dierential equation for x = x(t) given by
x = f(x). (7.1)
We say that x
is a xed point, or equilibrium point, of (7.1) if f(x
) = 0. At
the xed point, x = 0. The terminology xed point is used since the solution to
(7.1) with initial condition x(0) = x
is x(t) = x
for all time t.

A xed point, however, can be stable or unstable. A xed point is said to be
stable if a small perturbation of the solution from the xed point decays in time;
if the perturbation grows in time, the xed point is said to be unstable. We can
determine stability by a linear analysis. Let x = x
+ (t), where represents

a small perturbation of the solution from the xed point x
. Using x = since
x
is constant, f(x
) = 0 since x
is a xed point, and Taylor series expanding

about = 0, we have
= f(x
+)
= f(x
) +f
(x
) +. . .
= f
(x
) +. . . .
The omitted terms in the Taylor series expansion are proportional to
2
, and can
be made negligible with respect to the kept term, proportional to , by taking
77
78 CHAPTER 7. BIFURCATION THEORY
(0) suciently small. Therefore, at least over short times, the dierential
equation to be considered, = f
(x
), is linear and has by now the familiar

solution
(t) = (0)e
f
(x
)t
.
The perturbation of the xed point solution x(t) = x
thus decays exponentially

if f
(x
) < 0, and we say the xed point is stable. If f
(x
) > 0, the perturbation

grows exponenitally and we say the xed point is unstable. If f
(x
) = 0, we
say the xed point is marginally stable, and the next higher-order term in the
Taylor series expansion must be considered.
Example: Find all the xed points of the logistic equation x = x(1x)
and determine their stability.
There are two xed points at which x = 0, given by x
= 0 and x
= 1. Stability
of these equilibrium points may be determined by considering the derivative of
f(x) = x(1 x). We have f
(x) = 1 2x. Therefore, f
(0) = 1 > 0 so that

x
= 0 is an unstable xed point, and f
(1) = 1 < 0 so that x
= 1 is a stable
xed point. Indeed, we have previously found that all solutions approach the
stable xed point asymptotically.
The idea of xed points and stability can be extended to higher-order systems
of odes. Here, we consider a two-dimensional system and will need to make use
of the two-dimensional Taylor series expansion of a function F(x, y) about the
origin. In general, the Taylor series of F(x, y) is given by
F(x, y) = F +x
F
x
+y
F
y
+
1
2
_
x
2
2
F
x
2
+ 2xy

2
F
xy
+y
2
2
F
y
2
_
+. . . ,
where the function F and all of its partial derivatives on the right-hand-side are
evaluated at the origin. Note that the Taylor series is constructed so that all
partial derivatives of the left-hand-side match those of the right-hand-side at
the origin.
We now consider the two-dimensional system given by
x = f(x, y), y = g(x, y). (7.2)
The point (x
, y
) is said to be a xed point of (7.2) if f(x
, y
) = 0 and
g(x
, y
) = 0. Again, the local stability of a xed point can be determined by

a linear analysis. We let x(t) = x
+ (t) and y(t) = y
+ (t), where and

are small independent perturbations from the xed point. Making use of the
two dimensional Taylor series of f(x, y) and g(x, y) about the xed point, or
(, ) = (0, 0), we have
= f(x
+, y
+)
= f +
f
x
+
f
y
+. . .
=
f
x
+
f
y
+. . . .
= g(x
+, y
+)
= g +
g
x
+
g
y
+. . .
=
g
x
+
g
y
+. . . ,
7.1. FIXED POINTS AND STABILITY 79
0 0.5 1 1.5 2 2.5 3 3.5
0
0.5
1
1.5
2
2.5
x
y
phase space diagram
Figure 7.1: Phase space plot for two-dimensional nonlinear system.
where in the Taylor series f, g and all their partial derivatives are evaluated at
the xed point (x
, y
). Neglecting higher-order terms in the Taylor series, we

thus have a system of odes for the perturbation, given in matrix form as
d
dt
_

_
=
_
f
x
f
y
g
x
g
y
_
_

_
. (7.3)
The two-by-two matrix in (7.3) is called the Jacobian matrix at the xed point.
An eigenvalue analysis of the Jacobian matrix will typically yield two eigenvalues
1
and
2
. These eigenvalues may be real and distinct, complex conjugate pairs,
or repeated. The xed point is stable (all perturbations decay exponentially)
if both eigenvalues have negative real parts. The xed point is unstable (some
perturbations grow exponentially) if at least one of the eigenvalues has a positive
real part. Fixed points can be further classied as stable or unstable nodes,
unstable saddle points, or stable or unstable improper nodes.
Example: Find all the xed points of the system x = x(3 x 2y),
y = y(2 x y), and determine their stability.
view lecture
The xed points are determined by solving
f(x, y) = x(3 x 2y) = 0, g(x, y) = y(2 x y) = 0.
There are four xed points (x
, y
): (0, 0), (0, 2), (3, 0) and (1, 1). The Jacobian
matrix is given by
_
f
x
f
y
g
x
g
y
_
=
_
3 2x 2y 2x
y 2 x 2y
_
.
Stability of the xed points may be considered in turn. With J
the Jacobian
matrix evaluated at the xed point, we have
(x
, y
) = (0, 0) : J
=
_
3 0
0 2
_
.
The eigenvalues of J
are = 3, 2 so that the xed point (0, 0) is an unstable

node. Next,
(x
, y
) = (0, 2) : J
=
_
1 0
2 2
_
.
are = 1, 2 so that the xed point (0, 2) is a stable

node. Next,
(x
, y
) = (3, 0) : J
=
_
3 6
0 1
_
.
are = 3, 1 so that the xed point (3, 0) is also a

stable node. Finally,
(x
, y
) = (1, 1) : J
=
_
1 2
1 1
_
.
The characteristic equation of J
is given by (1 )
2
2 = 0, so that =
1
2. Since one eigenvalue is negative and the other positive the xed point
(1, 1) is an unstable saddle point. From our analysis of the xed points, one can
expect that all solutions will asymptote to one of the stable xed points (0, 2)
or (3, 0), depending on the initial conditions.
It is of interest to sketch the phase space diagram for this nonlinear system.
The eigenvectors associated with the unstable saddle point (1, 1) determine the
directions of the ow into and away from this xed point. The eigenvector
associated with the positive eigenvalue
1
= 1 +
2 can determined from the

rst equation of (J
1
I)v
1
= 0, or
2v
11
2v
12
= 0,
so that v
12
= (
2/2)v
11
. The eigenvector associated with the negative eigen-
value
1
= 1
2 satises v
22
= (
2/2)v
21
. The eigenvectors give the slope
of the lines with origin at the xed point for incoming (negative eigenvalue) and
outgoing (positive eigenvalue) trajectories. The outgoing trajectories have neg-
ative slope
2/2 and the incoming trajectories have positive slope

2/2. A
rough sketch of the phase space diagram can be made by hand (as demonstrated
in class). Here, a computer generated plot obtained from numerical solution of
the nonlinear coupled odes is presented in Fig. 7.1.
7.2 One-dimensional bifurcations
We now consider four classic bifurcations of one-dimensional nonlinear dier-
ential equations: saddle-node bifurcation, transcritical bifurcation, supercritical
pitchfork bifurcation, and subcritical pitchfork bifurcation. The corresponding
dierential equation will be written as
x = f
r
(x),
7.2. ONE-DIMENSIONAL BIFURCATIONS 81

r<0

r=0
r>0
dx/dt
x
(a)
(b)
x
*
r
Figure 7.2: Saddlenode bifurcation. (a) x versus x; (b) bifurcation diagram.
where the subscript r represents a parameter that results in a bifurcation when
varied across zero. The simplest dierential equations that exhibit these bifur-
cations are called the normal forms, and correspond to a local analysis (i.e.,
Taylor series expansion) of more general dierential equations around the xed
point, together with a possible rescaling of x.
7.2.1 Saddle-node bifurcation
view lecture
The saddle-node bifurcation results in xed points being created or destroyed.
The normal form for a saddle-node bifurcation is given by
x = r +x
2
.
The xed points are x
r. Clearly, two real xed points exist when

r < 0 and no real xed points exist when r > 0. The stability of the xed
points when r < 0 are determined by the derivative of f(x) = r + x
2
, given by
f
(x) = 2x. Therefore, the negative xed point is stable and the positive xed
point is unstable.
Graphically, we can illustrate this bifurcation in two ways. First, in Fig. 7.2(a),
we plot x versus x for the three parameter values corresponding to r < 0, r = 0
and r > 0. The values at which x = 0 correspond to the xed points, and arrows
are drawn indicating how the solution x(t) evolves (to the right if x > 0 and to
the left if x < 0. The stable xed point is indicated by a lled circle and the
unstable xed point by an open circle. Note that when r = 0, solutions converge
to the origin from the left, but diverge from the origin on the right. Second, in

r<0
dx/dt
x
(a)

r=0

r>0
x
*
r
(b)
Figure 7.3: Transcritical bifurcation. (a) x versus x; (b) bifurcation diagram.
Fig. 7.2(b), we plot a bifurcation diagram illustrating the xed point x
versus
the bifurcation parameter r. The stable xed point is denoted by a solid line
and the unstable xed points by a dashed line. Note that the two xed points
collide and annihilate at r = 0, and there are no xed points for r > 0.
7.2.2 Transcritical bifurcation
view lecture
A transcritical bifurcation occurs when there is an exchange of stabilities be-
tween two xed points. The normal form for a transcritical bifurcation is given
by
x = rx x
2
.
= 0 and x
= r. The derivative of the right-hand-side is

f
(x) = r 2x, so that f
(0) = r and f
(r) = r. Therefore, for r < 0, x
= 0
is stable and x
= r is unstable, while for r > 0, x
= r is stable and x
= 0 is
unstable. We see that there is an exchange of stabilities between the two xed
points as r passes through zero. The transcritical bifurcation is illustrated in
Fig. 7.3.
7.2.3 Supercritical pitchfork bifurcation
The pitchfork bifurcations occur in physical models where xed points appear
and disappear in pairs due to some intrinsic symmetry of the problem. The
normal form for the supercritical pitchfork bifurcation is given by
x = rx x
3
.

r<0
dx/dt
x
(a)

r=0

r>0
x
*
r
(b)
Figure 7.4: Supercritical pitchfork bifurcation. (a) x versus x; (b) bifurcation
diagram.
Note that the linear term results in exponential growth when r > 0 and the
nonlinear term stablizes this growth. The xed points are x
= 0 and x
r, the latter xed points existing only when r > 0. The derivative of f is
f
(x) = r 3x
2
so that f
(0) = r and f
r) = 2r. Therefore, the xed

point x
= 0 is stable for r < 0 and unstable for r > 0 while the xed points
x =
r exist and are stable for r > 0. Notice that the xed point x
= 0
becomes unstable as r crosses zero and two new stable xed points x
r
are born. The supercritical pitchfork bifurcation is illustrated in Fig. 7.4.
7.2.4 Subcritical pitchfork bifurcation
view lecture
In the subcritical case, the cubic term is destabilizing. The normal form (to
order x
3
) is
x = rx +x
3
.
= 0 and x
r, the latter xed points existing

only when r 0. The derivative of the right-hand-side is f
(x) = r + 3x
2
so
that f
(0) = r and f
r) = 2r. Therefore, the xed point x
= 0 is stable
for r < 0 and unstable for r > 0 while the xed points x =
r exist and are

unstable for r < 0. There are no stable xed points when r > 0.
The absence of stable xed points for r > 0 indicates that the neglect of
terms of higher-order in x than x
3
in the normal form may be unwarranted.
Keeping to the intrinsic symmetry of the equations (only odd powers of x) we
can add a stabilizing nonlinear term proportional to x
5
. The extended normal
x
*
r
Figure 7.5: Subcritical pitchfork bifurcation diagram.
form (to order x
5
) is
x = rx +x
3
x
5
,
and is somewhat more dicult to analyze. The xed points are solutions of
x(r +x
2
x
4
) = 0.
The xed point x
= 0 exists for all r, and four additional xed points can be

found from the solutions of the quadratic equation in x
2
:
x
=
_
1
2
(1
1 + 4r).
These xed points exist only if x
is real. Clearly, for the inner square-root to

be real, r 1/4. Also observe that 1
1 + 4r becomes negative for r > 0.

We thus have three intervals in r to consider, and these regions and their xed
points are
r <
1
4
: x
= 0 (one xed point);
1
4
< r < 0 : x
= 0, x
=
_
1
2
(1
1 + 4r) (ve xed points);

r > 0 : x
= 0, x
=
_
1
2
(1 +
1 + 4r) (three xed points).

Stability is determined from f
(x) = r +3x
2
5x
4
. Now, f
(0) = r so x
= 0 is
stable for r < 0 and unstable for r > 0. The calculation for the other four roots
can be simplied by noting that x
satises r + x
2
x
4
= 0, or x
4
= r + x
2
.
Therefore,
f
(x
) = r + 3x
2
5x
4
= r + 3x
2
5(r +x
2
)
= 4r 2x
2
= 2(2r +x
2
).
With x
2
=
1
2
(1
1 + 4r), we have
f
(x
) = 2
_
2r +
1
2
(1
1 + 4r)
_
=
_
(1 + 4r)
1 + 4r
_
=
1 + 4r
_
1 + 4r 1
_
.
Clearly, the plus root is always stable since f
(x
) < 0. The minus root exists

only for
1
4
< r < 0 and is unstable since f
(x
) > 0. We summarize the

stability of the various xed points:
r <
1
4
: x
= 0 (stable);
1
4
< r < 0 : x
= 0, (stable)
x
=
_
1
2
(1 +
1 + 4r) (stable);
x
=
_
1
2
(1
1 + 4r) (unstable);
r > 0 : x
= 0 (unstable)
x
=
_
1
2
(1 +
1 + 4r) (stable).
The bifurcation diagram is shown in Fig. 7.5. We can imagine what happens
to the solution to the ode as r increases from negative values, supposing there
is some noise in the system so that x(t) uctuates around a stable xed point.
For r < 1/4, the solution x(t) uctuates around the stable xed point x
= 0.
As r increases into the range 1/4 < r < 0, the solution will remain close to the
stable xed point x
= 0. However, a catastrophe occurs as as soon as r > 0.

The x
= 0 xed point is lost and the solution will jump up (or down) to the
only remaining xed point. A similar catastrophe can happen as r decreases
from positive values. In this case, the jump occurs as soon as r < 1/4. Since
the behavior of x(t) is dierent depending on whether we increase or decrease
r, we say that the system exhibits hysteresis. The existence of a subcritical
pitchfork bifurcation can be very dangerous in engineering applications since
a small change in a problems parameters can result in a large change in the
equilibrium state. Physically, this can respond to the collapse of a structure.
7.2.5 Application: a mathematical model of a shery
view lecture
We illustrate the utility of bifurcation theory by analyzing a simple model of a
shery. We utilize the logistic equation (see 2.3.8) to model a sh population
in the absence of shing. To model shing, we assume that the government has
established shing quotas so that at most a total of C sh per year may be
caught. We assume that when the sh population is at the carrying capacity
of the environment, sherman can catch nearly their full quota. When the sh
population drops to lower values, sh may be harder to nd and the catch rate
may fall below C, eventually going to zero as the sh population diminishes.
Combining the logistic equation together with a simple model of shing, we
propose the mathematical model
dN
dt
= rN
_
1
N
K
_
CN
A+N
, (7.4)
x
*
c
1
0
(1+)
2
/4
Figure 7.6: Fishery bifurcation diagram.
where N is the sh population size, t is time, r is the maximum potential growth
rate of the sh population, K is the carrying capacity of the environment, C is
the maximum rate at which sh can be caught, and A is a constant satisfying
A < K that is used to model the idea that sh become harder to catch when
scarce.
We nondimensionalize (7.4) using x = N/K, = rt, c = C/rK, = A/K:
dx
d
= x(1 x)
cx
+x
. (7.5)
Note that 0 x 1, c > 0 and 0 < < 1. We wish to qualitatively describe
the equilibrium solutions of (7.5) and the bifurcations that may occur as the
nondimensional catch rate c increases from zero. Practically, a government
would like to issue each year as large a catch quota as possible without adversely
aecting the number of sh that may be caught in subsequent years.
The xed points of (7.5) are x
= 0, valid for all c, and the solutions to

x
2
(1 )x + (c ) = 0, or
x
=
1
2
_
(1 )
_
(1 +)
2
4c
_
. (7.6)
The xed points given by (7.6) are real only if c <
1
4
(1 + )
2
. Furthermore,
x
> 0 only when (1 ) >

_
(1 +)
2
4c, or after some algebra, c > . We
therefore need to consider three intervals over which the following xed points
exist:
0 c : x
= 0, x
+
=
1
2
_
(1 ) +
_
(1 +)
2
4c
_
;
< c <
1
4
(1 +)
2
: x
= 0, x
=
1
2
_
(1 )
_
(1 +)
2
4c
_
;
7.3. TWO-DIMENSIONAL BIFURCATIONS 87
c >
1
4
(1 +)
2
: x
= 0.
The stability of the xed points can be determined with rigour analytically or
graphically. Here, we simply apply biological intuition together with knowledge
of the types of one dimensional bifurcations. An intuitive argument is made
simpler if we consider c decreasing from large values. When c is large, that is
c >
1
4
(1+)
2
, too many sh are being caught and our intuition suggests that the
sh population goes extinct. Therefore, in this interval, the single xed point
x
= 0 must be stable. As c decreases, a bifurcation occurs at c =

1
4
(1 + )
2
introducing two additional xed points at x
= (1 )/2. The type of one

dimensional bifurcation in which two xed points are created as a square root
becomes real is a saddlenode bifurcation, and one of the xed points will be
stable and the other unstable. Following these xed points as c 0, we observe
that the upper branch x
+
goes to one, which is the appropriate stable xed
point when there is no shing. We therefore conclude that the upper solution
x
+
is stable and the lower solution x
is unstable. As c decreases further from

this bifurcation, the lower xed point x
collides with the xed point x
= 0 at
c = . This appears to be a transcritical bifurcation and assuming an exchange
of stability occurs, we must have the xed point x
= 0 becoming unstable for

c < . The resulting bifurcation diagram is shown in Fig. 7.6.
The purpose of simple mathematical models applied to complex ecological
problems is to oer some insight. Here, we have learned that overshing (in
the model c >
1
4
(1 + )
2
) during one year can potentially result in a sudden
collapse of the sh catch in subsequent years, so that governments need to be
particularly cautious when contemplating increases in shing quotas.
7.3 Two-dimensional bifurcations
view lecture
All the one-dimensional bifurcations can also occur in two-dimensions along
one of the directions. In addition, a new type of bifurcation can also occur
in two-dimensions. Suppose there is some control parameter . Furthermore,
suppose that for < 0, a two-dimensional system approaches a xed point by
exponentially-damped oscillations. We know that the Jacobian matrix at the
xed point with < 0 will have complex conjugate eigenvalues with negative real
parts. Now suppose that when > 0 the real parts of the eigenvalues become
positive so that the xed point becomes unstable. This change in stability
of the xed point is called a Hopf bifurcation. The Hopf bifurcation comes
in two types: supercritical Hopf bifurcation and subcritical Hopf bifurcation.
For the supercritical Hopf bifurcation, as increases slightly above zero, the
resulting oscillation around the now unstable xed point is quickly stabilized at
small amplitude. This stable orbit is called a limit cycle. For the subcritical
Hopf bifurcation, as increases slightly above zero, the limit cycle immediately
jumps to large amplitude.
7.3.1 Supercritical Hopf bifurcation
A simple example of a supercritical Hopf bifurcation can be given in polar
coordinates:
r = r r
3
,
= +br
2
,
where x = r cos and y = r sin. The parameter controls the stability of the
xed point at the origin, the parameter is the frequency of oscillation near
the origin, and the parameter b determines the dependence of the oscillation
frequency at larger amplitude oscillations. Although we include b for generality,
our qualitative analysis of these equations will be independent of b.
The equation for the radius is of the form of the supercritical pitchfork
bifurcation. The xed points are r
= 0 and r
=

(note that r > 0), and
the former xed point is stable for < 0 and the latter is stable for > 0. The
transition of the eigenvalues of the Jacobian from negative real part to positive
real part can be seen if we transform these equations to cartesian coordinates.
We have using r
2
= x
2
+y
2
x = r cos

r sin
= (r r
3
) cos ( +br
2
)r sin
= x (x
2
+y
2
)x y b(x
2
+y
2
)y
= x y (x
2
+y
2
)(x +by);
y = r sin +

r cos
= (r r
3
) sin + ( +br
2
)r cos
= y (x
2
+y
2
)y +x +b(x
2
+y
2
)x
= x +y (x
2
+y
2
)(y bx).
The stability of the origin is determined by the Jacobian matrix evaluated at
the origin. The nonlinear terms in the equation vanish and the Jacobian matrix
at the origin is given by
J =
_

_
.
The eigenvalues are the solutions of ( )
2
+
2
= 0, or = i. As
increases from negative to positive values, exponentially damped oscillations
change into exponentially growing oscillations. The nonlinear terms in the equa-
tions stabilize the growing oscillations into a limit cycle.
7.3.2 Subcritical Hopf bifurcation
The analogous example of a subcritical Hopf bifurcation is given by
r = r +r
3
r
5
,
= +br
2
.
Here, the equation for the radius is of the form of the subcritical pitchfork
bifurcation. As increases from negative to positive values, the origin becomes
unstable and exponentially growing oscillations increase until the radius reaches
7.3. TWO-DIMENSIONAL BIFURCATIONS 89
a stable xed point far away from the origin. In practice, it may be dicult
to tell analytically if a Hopf bifurcation is supercritical or subcritical from the
equations of motion. Computational solution, however, can quickly distinguish
between the two types.
Chapter 8
Partial dierential
equations
Dierential equations containing partial derivatives with two or more inde-
pendent variables are called partial dierential equations (pdes). These equa-
tions are of fundamental scientic interest but are substantially more dicult
to solve, both analytically and computationally, than odes. In this chapter, we
will derive two fundamental pdes and show how to solve them.
8.1 Derivation of the diusion equation
view lecture
To derive the diusion equation in one spacial dimension, we imagine a still
liquid in a long pipe of constant cross sectional area. A small quantity of dye
is placed in a cross section of the pipe and allowed to diuse up and down the
pipe. The dye diuses from regions of higher concentration to regions of lower
concentration. We dene u(x, t) to be the concentration of the dye at position
x along the pipe, and we wish to nd the pde satised by u. It is useful to keep
track of the units of the various quantities involved in the derivation and we
introduce the bracket notation [X] to mean the units of X. Relevant dimensional
units used in the derivation of the diusion equation are mass, length, and time,
denoted as m, l, and t, respectively. Assuming that the dye concentration is
uniform in every cross section of the pipe, the dimensions of concentration used
here are [u] = m/l. The mass of dye in the innitesimal pipe volume located
between position x
1
and position x
2
at time t is given to order x = x
2
x
1
by
M = u(x, t)x.
The mass of dye in this innitesimal pipe volume changes by diusion into and
out of the cross sectional ends situated at position x
1
and x
2
(Figure 8.1).
We assume the rate of diusion is proportional to the concentration gradient,
a relationship known as Ficks law of diusion. Ficks law of diusion assumes
the mass ux J, with units [J] = m/t across a cross section of the pipe is given
by
J = Du
x
, (8.1)
91
92 CHAPTER 8. PARTIAL DIFFERENTIAL EQUATIONS
J (x
2
)
J (x
1
)
X
1
X
2
Figure 8.1: Derivation of the diusion equation.
where the diusion constant D > 0 has units [D] = l
2
/t. The mass ux is
opposite in sign to the gradient of concentration. The time rate of change in the
mass of dye between x
1
and x
2
is given by the dierence between the mass ux
into and the mass ux out of the innitesimal cross sectional volume. When
u
x
< 0, J > 0 and the mass ows into the volume at position x
1
and out of the
volume at position x
2
. On the other hand, when u
x
> 0, J < 0 and the mass
ows out of the volume at position x
1
and into the volume at position x
2
. In
both cases, the time rate of change of the dye mass is given by
dM
dt
= J(x
1
, t) J(x
2
, t),
or rewriting in terms of u(x, t):
u
t
(x, t)x = D(u
x
(x
2
, t) u
x
(x
1
, t)) .
Dividing by x and taking the limit x 0 results in the diusion equation:
u
t
= Du
xx
.
We note that the diusion equation is identical to the heat conduction equation,
where u is temperature, and the constant D (commonly written as ) is the
thermal conductivity.
8.2 Derivation of the wave equation
view lecture
To derive the wave equation in one spacial dimension, we imagine an elastic
string that undergoes small amplitude transverse vibrations. We dene u(x, t)
to be the vertical displacement of the string from the x-axis at position x and
time t, and we wish to nd the pde satised by u. We dene to be the
constant mass density of the string, T the tension of the string, and the angle
between the string and the horizonal line. We consider an innitesimal string
element located between x
1
and x
2
, with x = x
2
x
1
, as shown in Fig. 8.2.
The governing equations are Newtons law of motion for the horizontal and
8.2. DERIVATION OF THE WAVE EQUATION 93
Figure 8.2: Derivation of the wave equation.
vertical acceleration of our innitesimal string element, and we assume that the
string element only accelerates vertically. Therefore, the horizontal forces must
balance and we have
T
2
cos
2
= T
1
cos
1
.
The vertical forces result in a vertical acceleration, and with u
tt
the vertical
acceleration of the string element and
x
2
+ u
2
= x
_
1 +u
2
x
its mass,
where we have used u
x
= u/x, exact as x 0, we have
x
_
1 +u
2
x
u
tt
= T
2
sin
2
T
1
sin
1
.
We now make the assumption of small vibrations, that is u << x, or equiv-
alently u
x
<< 1. Note that [u] = l so that u
x
is dimensionless. With this
approximation, to leading-order in u
x
we have
cos
2
= cos
1
= 1,
sin
2
= u
x
(x
2
, t), sin
1
= u
x
(x
1
, t),
and
_
1 +u
2
x
= 1.
Therefore, to leading order T
1
= T
2
= T, (i.e., the tension in the string is
approximately constant), and
xu
tt
= T (u
x
(x
2
, t) u
x
(x
1
, t)) .
Dividing by x and taking the limit x 0 results in the wave equation
u
tt
= c
2
u
xx
,
where c
2
= T/. Since [T] = ml/t
2
and [] = m/l, we have [c
2
] = l
2
/t
2
so that
c has units of velocity.
8.3 Fourier Series
view lecture
view lecture
Our solution of the diusion and wave equations will require use of a Fourier
series. A periodic function f(x) with period 2L, can be represented as a Fourier
series in the form
f(x) =
a
0
2
+
n=1
_
a
n
cos
nx
L
+b
n
sin
nx
L
_
. (8.2)
Determination of the coecients a
0
, a
1
, a
2
, . . . and b
1
, b
2
, b
3
, . . . makes use of
orthogonality relations for sine and cosine. We rst dene the widely used
Kronecker delta
nm
as
nm
=
_
1 if n = m;
0 otherwise.
The orthogonality relations for n and m positive integers are then given with
compact notation as the integration formulas
_
L
L
cos
_
mx
L
_
cos
_
nx
L
_
dx = L
nm
, (8.3)
_
L
L
sin
_
mx
L
_
sin
_
nx
L
_
dx = L
nm
, (8.4)
_
L
L
cos
_
mx
L
_
sin
_
nx
L
_
dx = 0. (8.5)
As an example of the integration technique used to obtain these results, for n
and m positive integers and n = m, we derive (8.4) by
_
L
L
sin
_
mx
L
_
sin
_
nx
L
_
dx
=
1
2
_
L
L
_
cos
_
(mn)x
L
_
cos
_
(m+n)x
L
__
dx
=
L
2
_
1
mn
sin
_
(mn)x
L
_
1
m+n
sin
_
(m+n)x
L
__
L
L
= 0.
For m and n positive integers and m = n, however,
_
L
L
sin
2
_
nx
L
_
dx =
1
2
_
L
L
_
1 cos
2nx
L
_
dx
=
1
2
_
x
L
2n
sin
2nx
L
_
L
L
= L.
Integration formulas given by (8.3) and (8.5) can be similarly derived.
8.3. FOURIER SERIES 95
To determine the coecients a
0
, a
1
, a
2
, . . ., we multiply both sides of (8.2)
by cos (nx/L) with n a nonnegative integer (changing the dummy summation
variable from n to m) and integrate over x from L to L, assuming that the
integration can be done term by term in the innite sum:
_
L
L
f(x) cos
nx
L
dx =
a
0
2
_
L
L
cos
nx
L
dx
+
m=1
_
a
m
_
L
L
cos
nx
L
cos
mx
L
dx +b
m
_
L
L
cos
nx
L
sin
mx
L
dx
_
.
If n = 0, then the second and third integrals on the right-hand-side are zero
because the cosine and sine have period 2L, the integrand in the rst integral on
the right-hand-side is a constant, and the right-hand-side becomes La
0
. If n is
a positive integer, then the rst integral on the right-hand-side is zero because
the cosine has period 2L and the third integral on the right-hand-side is zero
because of the orthogonality of cosine and sine, (8.5). The second integral on
the right-hand-side is performed using (8.3) to obtain
_
L
L
f(x) cos
nx
L
dx =
m=1
La
m
nm
= La
n
,
where all the terms in the summation except m = n are zero by virtue of the
Kronecker delta. We therefore obtain for n = 0, 1, 2, . . .
a
n
=
1
L
_
L
L
f(x) cos
nx
L
dx. (8.6)
To determine the coecients b
1
, b
2
, b
3
, . . ., we multiply both sides of (8.2) by
sin(nx/L) with n a positive integer (again changing the dummy summation
variable from n to m) and integrate:
_
L
L
f(x) sin
nx
L
dx =
a
0
2
_
L
L
sin
nx
L
dx
+
m=1
_
a
m
_
L
L
sin
nx
L
cos
mx
L
dx +b
m
_
L
L
sin
nx
L
sin
mx
L
dx
_
.
The rst integral on the right-hand-side is zero because of the periodicity of
sine and the second integral on the right-hand-side is zero because of the or-
thogonality of cosine and sine, (8.5). The third integral on the right-hand-side
is performed using (8.4) to obtain
_
L
L
f(x) sin
nx
L
dx =
m=1
Lb
m
nm
= Lb
n
,
so that for n = 1, 2, 3, . . .
b
n
=
1
L
_
L
L
f(x) sin
nx
L
dx. (8.7)
Our results for the Fourier series of a function f(x) with period 2L are thus
given by (8.2), (8.6) and (8.7).
The Fourier series simplies if f(x) is an even function such that f(x) =
f(x), or an odd function such that f(x) = f(x). Use will be made of the
facts that cos (nx/L) is an even function and sin(nx/L) is an odd function;
the product of two even functions is an even function; the product of two odd
functions is an even function; the product of an even and an odd function is an
odd function, and; that depending on whether g(x) is an even or odd function,
_
L
L
g(x)dx = 2
_
L
0
g(x)dx, g(x) even;
_
L
L
g(x)dx = 0, g(x) odd.
We examine the Fourier series for an even or odd function in turn. First, if f(x)
is even, then from (8.6) and (8.7) and our facts about even and odd functions,
a
n
=
2
L
_
L
0
f(x) cos
nx
L
dx, (8.8)
b
n
= 0.
The Fourier series for an even function with period 2L is thus given by the
Fourier cosine series
f(x) =
a
0
2
+
n=1
a
n
cos
nx
L
, f(x) even. (8.9)
Second, if f(x) is odd, then
a
n
= 0,
b
n
=
2
L
_
L
0
f(x) sin
nx
L
dx; (8.10)
and the Fourier series for an odd function with period 2L is given by the Fourier
sine series
f(x) =
n=1
b
n
sin
nx
L
, f(x) odd. (8.11)
With Fourier series now included in our applied mathematics toolbox, we are
ready to solve the diusion and wave equations in bounded domains.
8.4 Solution of the diusion equation
8.4.1 Homogeneous boundary conditions
view lecture
view lecture
We consider one dimensional diusion in a pipe of length L, and solve the
diusion equation for the concentration u(x, t),
u
t
= Du
xx
, 0 x L, t > 0. (8.12)
8.4. SOLUTION OF THE DIFFUSION EQUATION 97
Both initial and boundary conditions are required for a unique solution. That
is, we assume the initial concentration distribution in the pipe is given by
u(x, 0) = f(x), 0 x L. (8.13)
Furthermore, we assume that boundary conditions are given at the ends of
the pipes. When the concentration value is specied at the boundaries, the
boundary conditions are called Dirichlet boundary conditions. As the simplest
example, we assume here homogeneous Dirichlet boundary conditions, that is
zero concentration of dye at the ends of the pipe, which could occur if the ends
of the pipe open up into large reservoirs of clear solution,
u(0, t) = 0, u(L, t) = 0, t > 0. (8.14)
We will later also discuss inhomogeneous Dirichlet boundary conditions and
homogeneous Neumann boundary conditions, for which the derivative of the
concentration is specied to be zero at the boundaries. Note that if f(x) is iden-
tically zero, then the trivial solution u(x, t) = 0 saties the dierential equation
and the initial and boundary conditions and is therefore the unique solution of
the problem. In what follows, we will assume that f(x) is not identically zero
so that we need to nd a solution dierent than the trivial solution.
The solution method we use is called separation of variables. We assume
that u(x, t) can be written as a product of two other functions, one dependent
only on position x and the other dependent only on time t. That is, we make
the ansatz
u(x, t) = X(x)T(t). (8.15)
Whether this ansatz will succeed depends on whether the solution indeed has
this form. Substituting (8.15) into (8.12), we obtain
XT
= DX
T,
which we rewrite by separating the x and t dependence to opposite sides of the
equation:
X
X
=
1
D
T
T
.
The left hand side of this equation is independent of t and the right hand side is
independent of x. Both sides of this equation are therefore independent of both
x and t and equal to a constant. Introducing as the separation constant, we
have
X
X
=
1
D
T
T
= ,
and we obtain the two ordinary dierential equations
X
+X = 0, T
+DT = 0. (8.16)
Because of the boundary conditions, we must rst consider the equation for
X(x). To solve, we need to determine the boundary conditions at x = 0 and
x = L. Now, from (8.14) and (8.15),
u(0, t) = X(0)T(t) = 0, t > 0.
Since T(t) is not identically zero for all t (which would result in the trivial
solution for u), we must have X(0) = 0. Similarly, the boundary condition at
x = L requires X(L) = 0. We therefore consider the two-point boundary value
problem
X
+X = 0, X(0) = X(L) = 0. (8.17)

The equation given by (8.17) is called an ode eigenvalue problem. The allowed
values of and the corresponding functions X(x) are called the eigenvalues and
eigenfunctions of the dierential equation. Since the form of the general solution
of the ode depends on the sign of , we consider in turn the cases > 0, < 0
and = 0. For > 0, we write =
2
and determine the general solution of
X
+
2
X = 0
to be
X(x) = Acos x +Bsinx.
Applying the boundary condition at x = 0, we nd A = 0. The boundary
condition at x = L then yields
BsinL = 0.
The solution B = 0 results in the trivial solution for u and can be ruled out.
Therefore, we must have
sinL = 0,
which is an equation for the eigenvalue . The solutions are
= n/L,
where n is an integer. We have thus determined the eigenvalues =
2
> 0 to
be
n
= (n/L)
2
, n = 1, 2, 3, . . . , (8.18)
with corresponding eigenfunctions
X
n
= sin(nx/L). (8.19)
For < 0, we write =
2
and determine the general solution of
X
2
X = 0
to be
X(x) = Acoshx +Bsinhx,
where we have previously introduced the hyperbolic sine and cosine functions
in 3.2.1. Applying the boundary condition at x = 0, we nd A = 0. The
boundary condition at x = L then yields
BsinhL = 0,
which for = 0 has only the solution B = 0. Therefore, there is no nontrivial
solution for u with < 0. Finally, for = 0, we have
X
= 0,
with general solution
X(x) = A+Bx.
The boundary condition at x = 0 and x = L yields A = B = 0 so again there is
no nontrivial solution for u with = 0.
We now turn to the equation for T(t). The equation corresponding to the
eigenvalue
n
, using (8.18), is given by
T
+
_
n
2
2
D/L
2
_
T = 0,
which has solution proportional to
T
n
= e
n
2
2
Dt/L
2
. (8.20)
Therefore, multiplying the solutions given by (8.19) and (8.20), we conclude
that the functions
u
n
(x, t) = sin(nx/L)e
n
2
2
Dt/L
2
(8.21)
satisfy the pde given by (8.12) and the boundary conditions given by (8.14) for
every positive integer n.
The principle of linear superposition for homogeneous linear dierential
equations then states that the general solution to (8.12) and (8.14) is given
by
u(x, t) =
n=1
b
n
u
n
(x, t)
=
n=1
b
n
sin(nx/L)e
n
2
2
Dt/L
2
. (8.22)
The nal solution step is to satisfy the initial conditions given by (8.13). At
t = 0, we have
f(x) =
n=1
b
n
sin(nx/L). (8.23)
We immediately recognize (8.23) as a Fourier sine series (8.11) for an odd func-
tion f(x) with period 2L. Equation (8.23) is a periodic extension of our original
f(x) dened on 0 x L, and is an odd function because of the boundary
condition f(0) = 0. From our solution for the coecients of a Fourier sine series
(8.10), we determine
b
n
=
2
L
_
L
0
f(x) sin
nx
L
dx. (8.24)
Thus the solution to the diusion equation with homogeneous Dirichlet bound-
ary conditions dened by (8.12), (8.13) and (8.14) is given by (8.22) with the
b
n
coecients computed from (8.24).
Example: Determine the concentration of a dye in a pipe of length L,
where the dye has unit mass and is initially concentrated at the center
of the pipe, and the ends of the pipe are held at zero concentration
The governing equation for concentration is the diusion equation. We model
the initial concentration of the dye by a delta-function centered at x = L/2,
that is, u(x, 0) = f(x) = (x L/2). Therefore, from (8.24),
b
n
=
2
L
_
L
0
(x
L
2
) sin
nx
L
dx
=
2
L
sin(n/2)
=
_
_
_
2/L if n = 1, 5, 9, . . .;
2/L if n = 3, 7, 11, . . .;
0 if n = 2, 4, 6, . . ..
With b
n
determined, the solution for u(x, t) given by (8.22) can be written as
u(x, t) =
2
L
n=0
(1)
n
sin
_
(2n + 1)x
L
_
e
(2n+1)
2
2
Dt/L
2
.
When t >> L
2
/D, the leading-order term in the series is a good approximation
and is given by
u(x, t)
2
L
sin(x/L)e
2
Dt/L
2
.
8.4.2 Inhomogeneous boundary conditions
view lecture
Consider a diusion problem where one end of the pipe has dye of concentration
held constant at C
1
and the other held constant at C
2
, which could occur if the
ends of the pipe had large reservoirs of uid with dierent concentrations of
dye. With u(x, t) the concentration of dye, the boundary conditions are given
by
u(0, t) = C
1
, u(L, t) = C
2
, t > 0.
The concentration u(x, t) satises the diusion equation with diusivity D:
u
t
= Du
xx
.
If we try to solve this problem directly using separation of variables, we will
run into trouble. Applying the inhomogeneous boundary condition at x = 0
directly to the ansatz u(x, t) = X(x)T(t) results in
u(0, t) = X(0)T(t) = C
1
;
so that
X(0) = C
1
/T(t).
However, our separation of variables ansatz assumes X(x) to be independent of
t! We therefore say that inhomogeneous boundary conditions are not separable.
The proper way to solve a problem with inhomogeneous boundary conditions
is to transform it into another problem with homogeneous boundary conditions.
As t , we assume that a stationary concentration distribution v(x) will
attain, independent of t. Since v(x) must satisfy the diusion equation, we have
v
(x) = 0, 0 x L,
with general solution
v(x) = A+Bx.
Since v(x) must satisfy the same boundary conditions of u(x, t), we have v(0) =
C
1
and v(L) = C
2
, and we determine A = C
1
and B = (C
2
C
1
)/L.
We now express u(x, t) as the sum of the known asymptotic stationary con-
centration distribution v(x) and an unknown transient concentration distribu-
tion w(x, t):
u(x, t) = v(x) +w(x, t).
Substituting into the diusion equation, we obtain
t
(v(x) +w(x, t)) = D

2
x
2
(v(x) +w(x, t))
or
w
t
= Dw
xx
,
since v
t
= 0 and v
xx
= 0. The boundary conditions satised by w are
w(0, t) = u(0, t) v(0) = 0,
w(L, t) = u(L, t) v(L) = 0,
so that w is observed to satisfy homogeneous boundary conditions. If the initial
conditions are given by u(x, 0) = f(x), then the initial conditions for w are
w(x, 0) = u(x, 0) v(x)
= f(x) v(x).
The resulting equations may then be solved for w(x, t) using the technique for
homogeneous boundary conditions, and u(x, t) subsequently determined.
8.4.3 Pipe with closed ends
view lecture
view lecture
There is no diusion of dye through the ends of a sealed pipe. Accordingly,
the mass ux of dye through the pipe ends, given by (8.1), is zero so that the
boundary conditions on the dye concentration u(x, t) becomes
u
x
(0, t) = 0, u
x
(L, t) = 0, t > 0, (8.25)
which are known as homogeneous Neumann boundary conditions. Again, we
apply the method of separation of variables and as before, we obtain the two
ordinary dierential equations given by (8.16). Considering rst the equation
for X(x), the appropriate boundary conditions are now on the rst derivative
of X(x), and we must solve
X
+X = 0, X
(0) = X
(L) = 0. (8.26)
Again, we consider in turn the cases > 0, < 0 and = 0. For > 0, we
write =
2
and determine the general solution of (8.26) to be
X(x) = Acos x +Bsinx,
so that taking the derivative
X
(x) = Asinx +Bcos x.

Applying the boundary condition X
(0) = 0, we nd B = 0. The boundary

condition at x = L then yields
AsinL = 0.
The solution A = 0 results in the trivial solution for u and can be ruled out.
Therefore, we must have
sinL = 0,
with solutions
= n/L,
where n is an integer. We have thus determined the eigenvalues =
2
> 0 to
be
n
= (n/L)
2
, n = 1, 2, 3, . . . , (8.27)
X
n
= cos (nx/L). (8.28)
For < 0, we write =
2
and determine the general solution of (8.26) to be
X(x) = Acoshx +Bsinhx,
X
(x) = Asinhx +Bcoshx.

Applying the boundary condition X
(0) = 0 yields B = 0. The boundary

condition X
(L) = 0 then yields

AsinhL = 0,
which for = 0 has only the solution A = 0. Therefore, there is no nontrivial
solution for u with < 0. Finally, for = 0, the general solution of (8.26) is
X(x) = A+Bx,
X
(x) = B.
The boundary condition X
(0) = 0 yields B = 0; X
(L) = 0 is then trivially

satised. Therefore, we have an additional eigenvalue and eigenfunction given
by
0
= 0, X
0
(x) = 1,
which can be seen as extending the formula obtained for eigenvalues and eigen-
vectors for positive given by (8.27) and (8.28) to n = 0.
We now turn to the equation for T(t). The equation corresponding to eigen-
value
n
, using (8.27), is given by
T
+
_
n
2
2
D/L
2
_
T = 0,
which has solution proportional to
T
n
= e
n
2
2
Dt/L
2
, (8.29)
valid for n = 0, 1, 2, . . . . Therefore, multiplying the solutions given by (8.28)
and (8.29), we conclude that the functions
u
n
(x, t) = cos (nx/L)e
n
2
2
Dt/L
2
(8.30)
satisfy the pde given by (8.12) and the boundary conditions given by (8.25) for
every nonnegative integer n.
The principle of linear superposition then yields the general solution as
u(x, t) =
n=0
c
n
u
n
(x, t)
=
a
0
2
+
n=1
a
n
cos (nx/L)e
n
2
2
Dt/L
2
, (8.31)
where we have redened the constants so that c
0
= a
0
/2 and c
n
= a
n
, n =
1, 2, 3, . . .. The nal solution step is to satisfy the initial conditions given by
(8.13). At t = 0, we have
f(x) =
a
0
2
+
n=1
a
n
cos (nx/L), (8.32)
which we recognize as a Fourier cosine series (8.9) for an even function f(x)
with period 2L. We have obtained a cosine series for the periodic extension of
f(x) because of the boundary condition f
(0) = 0, which is satised by an even

function. From our solution (8.8) for the coecients of a Fourier cosine series,
we determine
a
n
=
2
L
_
L
0
f(x) cos
nx
L
dx. (8.33)
Thus the solution to the diusion equation with homogenous Neumann bound-
ary conditions dened by (8.12), (8.13) and (8.25) is given by (8.31) with the
coecients computed from (8.33).
Example: Determine the concentration of a dye in a pipe of length
L, where the dye has unit mass and is initially concentrated at the
center of the pipe, and the ends of the pipe are sealed
Again we model the initial concentration of the dye by a delta-function centered
at x = L/2. From (8.33),
a
n
=
2
L
_
L
0
(x
L
2
) cos
nx
L
dx
=
2
L
cos (n/2)
=
_
_
_
2/L if n = 0, 4, 8, . . .;
2/L if n = 2, 6, 10, . . .;
0 if n = 1, 3, 5, . . ..
The rst two terms in the series for u(x, t) are given by
u(x, t) =
2
L
_
1/2 cos (2x/L)e
4
2
Dt/L
2
+. . .
_
.
Notice that as t , u(x, t) 1/L: the dye mass is conserved in the pipe
(since the pipe ends are sealed) and eventually diused uniformly throughout
the pipe of length L.
8.5 Solution of the wave equation
view lecture
8.5.1 Plucked string
view lecture
We assume an elastic string with xed ends is plucked like a guitar string. The
governing equation for u(x, t), the position of the string from its equilibrium
position, is the wave equation
u
tt
= c
2
u
xx
, (8.34)
with boundary conditions at the string ends located at x = 0 and L given by
u(0, t) = 0, u(L, t) = 0. (8.35)
Since the wave equation is second-order in time, initial conditions are required
for both the displacement of the string due to the plucking and the initial velocity
of the displacement. We assume
u(x, 0) = f(x), u
t
(x, 0) = 0, 0 x L. (8.36)
Again, we use the method of separation of variables and try the ansatz
u(x, t) = X(x)T(t). (8.37)
Substitution of our ansatz (8.37) into the wave equation (8.34) and separating
variables results in
X
X
=
1
c
2
T
T
= ,
yielding the two ordinary dierential equations
X
+X = 0, T
+c
2
T = 0. (8.38)
Again, we solve rst the equation for X(x). The appropriate boundary condi-
tions for X are given by
X(0) = 0, X(L) = 0, (8.39)
and we have solved this equation for X(x) previously in 8.4.1 (see (8.17)).
A nontrivial solution exists only when > 0, and our previously determined
solution was
n
= (n/L)
2
, n = 1, 2, 3, . . . , (8.40)
X
n
= sin(nx/L). (8.41)
8.5. SOLUTION OF THE WAVE EQUATION 105
With
n
specied, the T equation then becomes
T
n
+
n
2
2
c
2
L
2
T
n
= 0,
with general solution given by
T
n
(t) = Acos
nct
L
+Bsin
nct
L
. (8.42)
The second of the initial conditions given by (8.36) implies
u
t
(x, 0) = X(x)T
(0) = 0,
which can be satised only if T
(0) = 0. Applying this boundary condition to

(8.42), we nd B = 0. Combining our solution for X
n
(x), (8.41), and T
n
(t), we
have determined that
u
n
(x, t) = sin
nx
L
cos
nct
L
, n = 1, 2, 3, . . .
satises the wave equation, the boundary conditions at the string ends, and the
assumption of zero initial string velocity. Linear superposition of these solutions
results in the general solution for u(x, t) of the form
u(x, t) =
n=1
b
n
sin
nx
L
cos
nct
L
. (8.43)
The remaining condition to satisfy is the initial displacement of the string, the
rst equation of (8.36). We have
f(x) =
n=1
b
n
sin(nx/L),
which is observed to be a Fourier sine series (8.11) for an odd function with
period 2L. Therefore, the coecients b
n
are given by (8.10),
b
n
=
2
L
_
L
0
f(x) sin
nx
L
dx, n = 1, 2, 3, . . . . (8.44)
Our solution to the wave equation with plucked string is thus given by (8.43)
and (8.44). Notice that the solution is time periodic with period 2L/c. The
corresponding fundamental frequency is the reciprical of the period and is given
by f = c/2L. From our derivation of the wave equation in 8.2, the velocity c
is related to the density of the string and tension of the string T by c
2
= T/.
Therefore, the fundamental frequency (pitch) of our guitar string increases (is
raised) with increasing tension, decreasing string density, and decreasing string
length. Indeed, these are exactly the parameters used to construct, tune and
play a guitar.
The wave nature of our solution and the physical signicance of the velocity
c can be made more transparent if we make use of the trigonometric identity
sinxcos y =
1
2
(sin(x +y) + sin(x y)) .
With this identity, our solution (8.43) can be rewritten as
u(x, t) =
1
2
n=1
b
n
_
sin
n(x +ct)
L
+ sin
n(x ct)
L
_
. (8.45)
The rst and second sine functions can be interpreted as a travelling wave mov-
ing to the left or the right with velocity c. This can be seen by incrementing
time, t t + , and observing that the value of the rst sine function is un-
changed provided the position is shifted by x x c, and the second sine
function is unchanged provided x x + c. Two waves travelling in opposite
directions with equal amplitude results in a standing wave.
8.5.2 Hammered string
view lecture
In contrast to a guitar string that is plucked, a piano string is hammered. The
appropriate initial conditions for a piano string would be
u(x, 0) = 0, u
t
(x, 0) = g(x), 0 x L. (8.46)
Our solution proceeds as previously, except that now the homogeneous initial
condition on T(t) is T(0) = 0, so that A = 0 in (8.42). The wave equation
solution is therefore
u(x, t) =
n=1
b
n
sin
nx
L
sin
nct
L
. (8.47)
Imposition of initial conditions then yields
g(x) =
c
L
n=1
nb
n
sin
nx
L
.
The coecient of the Fourier sine series for g(x) is seen to be ncb
n
/L, and we
have
ncb
n
L
=
2
L
_
L
0
g(x) sin
nx
L
dx,
or
b
n
=
2
nc
_
L
0
g(x) sin
nx
L
dx.
8.5.3 General initial conditions
view lecture
If the initial conditions on u(x, t) are generalized to
u(x, 0) = f(x), u
t
(x, 0) = g(x), 0 x L, (8.48)
then the solution to the wave equation can be determined using the principle of
linear superposition. Suppose v(x, t) is the solution to the wave equation with
initial condition (8.36) and w(x, t) is the solution to the wave equation with
initial conditions (8.46). Then we have
u(x, t) = v(x, t) +w(x, t),
since u(x, t) satises the wave equation, the boundary conditions, and the initial
conditions given by (8.48).
8.6. THE LAPLACE EQUATION 107
Figure 8.3: Dirichlet problem for the Laplace equation in a rectangle.
8.6 The Laplace equation
The diusion equation in two spatial dimensions is
u
t
= D(u
xx
+u
yy
).
The steady-state solution, approached aymptotically in time, has u
t
= 0 so
that the steady-state solution u = u(x, y) satises the two-dimensional Laplace
equation
u
xx
+u
yy
= 0. (8.49)
We will consider the mathematical problem of solving the two dimensional
Laplace equation inside a rectangular or a circular boundary. The value of
u(x, y) will be specied on the boundaries, dening this problem to be of Dirich-
let type.
8.6.1 Dirichlet problem for a rectangle
view lecture
We consider the Laplace equation (8.49) for the interior of a rectangle 0 < x < a,
0 < y < b, (see Fig. 8.3), with boundary conditions
u(x, 0) = 0, u(x, b) = 0, 0 < x < a;
u(0, y) = 0, u(a, y) = f(y), 0 y b.
More general boundary conditions can be solved by linear superposition of so-
lutions.
We take our usual ansatz
u(x, y) = X(x)Y (y),
and nd after substitution into (8.49),
X
X
=
Y
Y
= ,
with the separation constant. We thus obtain the two ordinary dierential
equations
X
X = 0, Y
+Y = 0.
The homogeneous boundary conditions are X(0) = 0, Y (0) = 0 and Y (b) = 0.
We have already solved the equation for Y (y) in 8.4.1, and the solution yields
the eigenvalues
n
=
_
n
b
_
2
, n = 1, 2, 3, . . . ,
Y
n
(y) = sin
ny
b
.
The remaining X equation and homogeneous boundary condition is therefore
X
n
2
2
b
2
X = 0, X(0) = 0,
and the solution is the hyperbolic sine function
X
n
(x) = sinh
nx
b
,
times a constant. Writing u
n
= X
n
Y
n
, multiplying by a constant and summing
over n, yields the general solution
u(x, t) =
n=0
c
n
sinh
nx
b
sin
ny
b
.
The remaining inhomogeneous boundary condition u(a, y) = f(y) results in
f(y) =
n=0
c
n
sinh
na
b
sin
ny
b
,
which we recognize as a Fourier sine series for an odd function with period 2b,
and coecient c
n
sinh(na/b). The solution for the coecient is given by
c
n
sinh
na
b
=
2
b
_
b
0
f(y) sin
ny
b
dx.
8.6.2 Dirichlet problem for a circle
view lecture
view lecture
The Laplace equation is commonly written symbolically as
2
u = 0, (8.50)
where
2
is called the Laplacian, sometimes denoted as . The Laplacian can
be written in various coordinate systems, and the choice of coordinate systems
usually depends on the geometry of the boundaries. Indeed, the Laplace equa-
tion is known to be separable in 13 dierent coordinate systems! We have solved
the Laplace equation in two dimensions, with boundary conditions specied on
a rectangle, using
2
=

2
x
2
+

2
y
2
.
Here we consider boundary conditions specied on a circle, and write the Lapla-
cian in polar coordinates by changing variables from cartesian coordinates. Polar
coordinates is dened by the transformation (r, ) (x, y):
x = r cos , y = r sin; (8.51)
and the chain rule gives for the partial derivatives
u
r
=
u
x
x
r
+
u
y
y
r
,
u
=
u
x
x
+
u
y
y
. (8.52)
After taking the partial derivatives of x and y using (8.51), we can write the
transformation (8.52) in matrix form as
_
u/r
u/
_
=
_
cos sin
r sin r cos
__
u/x
u/y
_
. (8.53)
Inversion of (8.53) can be determined from the following result, commonly
proved in a linear algebra class. If
A =
_
a b
c d
_
, det A = 0,
then
A
1
=
1
det A
_
d b
c a
_
.
Therefore, since the determinant of the 2 2 matrix in (8.53) is r, we have
_
u/x
u/y
_
=
_
cos sin/r
sin cos /r
__
u/r
u/
_
. (8.54)
Rewriting (8.54) in operator form, we have
x
= cos

r

sin
r
,

y
= sin

r
+
cos
r
. (8.55)
To nd the Laplacian in polar coordinates with minimum algebra, we combine
(8.55) using complex variables as
x
+i

y
= e
i
_

r
+
i
r
_
, (8.56)
so that the Laplacian may be found by multiplying both sides of (8.56) by its
complex conjugate, taking care with the computation of the derivatives on the
right-hand-side:
2
x
2
+

2
y
2
= e
i
_

r
+
i
r
_
e
i
_

r

i
r
_
=

2
r
2
+
1
r
r
+
1
r
2
2
.
We have therefore determined that the Laplacian in polar coordinates is given
by
2
=

2
r
2
+
1
r
r
+
1
r
2
2
. (8.57)
We now consider the solution of the Laplace equation in a circle with radius
r < a subject to the boundary condition
u(a, ) = f(), 0 2. (8.58)
An additional boundary condition due to the use of polar coordinates is that
u(r, ) is periodic in with period 2. Furthermore, we will also assume that
u(r, ) is nite within the circle.
The method of separation of variables takes as our ansatz
u(r, ) = R(r)(),
and substitution into the Laplace equation (8.50) using (8.57) yields
R
+
1
r
R
+
1
r
2
R
= 0,
or
r
2
R
R
+r
R
R
=
= ,
where is the separation constant. We thus obtain the two ordinary dierential
equations
r
2
R
+rR
R = 0,
+ = 0.
The equation is solved assuming periodic boundary conditions with period
2. If < 0, then no periodic solution exists. If = 0, then can be constant.
If =
2
> 0, then
() = Acos +Bsin,
and the requirement that is periodic with period 2 forces to be an integer.
Therefore,
n
= n
2
, n = 0, 1, 2, . . . ,
n
() = A
n
cos n +B
n
sinn.
The R equation for each eigenvalue
n
then becomes
r
2
R
+rR
n
2
R = 0, (8.59)
which is an Euler equation. With the ansatz R = r
s
, (8.59) reduces to the
algebraic equation s(s 1) + s n
2
= 0, or s
2
= n
2
. Therefore, s = n, and
there are two real solutions when n > 0 and degenerate solutions when n = 0.
When n > 0, the solution for R(r) is
R
n
(r) = Ar
n
+Br
n
.
The requirement that u(r, ) is nite in the circle forces B = 0 since r
n
becomes
unbounded as r 0. When n = 0, the solution for R(r) is
R
n
(r) = A+Blnr,
and again nite u in the circle forces B = 0. Therefore, the solution for n =
0, 1, 2, . . . is R
n
= r
n
. Thus the general solution for u(r, ) may be written as
u(r, ) =
A
0
2
+
n=1
r
n
(A
n
cos n +B
n
sinn), (8.60)
where we have separated out the n = 0 solution to write our solution in a form
similar to the standard Fourier series given by (8.2). The remaining boundary
condition (8.58) species the values of u on the circle of radius a, and imposition
of this boundary condition results in
f() =
A
0
2
+
n=1
a
n
(A
n
cos n +B
n
sinn). (8.61)
Equation (8.61) is a Fourier series for the periodic function f() with period 2,
i.e., L = in (8.2). The Fourier coecients a
n
A
n
and a
n
B
n
are therefore given
by (8.6) and (8.7) to be
a
n
A
n
=
1
_
2
0
f() cos nd, n = 0, 1, 2, . . . ;
a
n
B
n
=
1
_
2
0
f() sinnd, n = 1, 2, 3, . . . . (8.62)
A remarkable fact is that the innite series solution for u(r, ) can be summed
explicitly. Substituting (8.62) into (8.60), we obtain
u(r, ) =
1
2
_
2
0
df()
_
1 + 2
n=1
_
r
a
_
n
(cos n cos n + sinn sinn)
_
=
1
2
_
2
0
df()
_
1 + 2
n=1
_
r
a
_
n
cos n( )
_
.
We can sum the innite series by writing 2 cos n( ) = e
in()
+e
in()
,
and using the sum of the geometric series

n=1
z
n
= z/(1 z) to obtain
1 + 2
n=1
_
r
a
_
n
cos n( ) = 1 +
n=1
_
re
i()
a
_n
+
n=1
_
re
i()
a
_n
= 1 +
_
re
i()
a re
i()
+ c.c.
_
=
a
2
r
2
a
2
2ar cos ( ) +r
2
.
Therefore,
u(r, ) =
a
2
r
2
2
_
2
0
f()
a
2
2ar cos ( ) +r
2
d,
an integral result for u(r, ) known as Poissons formula. As a trivial example,
consider the solution for u(r, ) if f() = F, a constant. Clearly, u(r, ) = F
satises both the Laplace equation and the boundary conditions so must be the
solution. You can verify that u(r, ) = F is indeed the solution by showing that
_
2
0
d
a
2
2ar cos ( ) +r
2
=
2
a
2
r
2
.

Introduction To Diff Erential Equations Jeff Rey R. Chasnov

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Introduction To Diff Erential Equations Jeff Rey R. Chasnov

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Introduction to Dierential

(x) or dy/dx, is dened

(x) = f(x). We say that F(x) is an antiderivative of

(x)dx = f(g(x)) +c.

(a), all other terms vanishing, etc. Com-

2gR 40 000 km/hr.

= 1. Therefore, the general solution for x is

4ac = 0), given by

) < 0. Therefore, both linearly independent homogeneous solutions

= t c. The translation of f(t) a

15(t 5)/4). (4.6)

(0) = 1. The numerical solutions, obtained using MATLAB,

+y = 0 has oscillatory sine and cosine solutions and y

and y(x) = Y (),

> 0. We need only consider x > 0 here and subsequently

and = lnx. By the chain

the real roots of (5.9), the general solution is

(Acos (ln|x|) +Bsin(ln|x|)) .

y = 0 for 0 x 1 with two-point

(0) (which is innite).

is a xed point, or equilibrium point, of (7.1) if f(x

for all time t.

+ (t), where represents

is a xed point, and Taylor series expanding

), is linear and has by now the familiar

thus decays exponentially

) < 0, and we say the xed point is stable. If f

) > 0, the perturbation

(x) = 1 2x. Therefore, f

(0) = 1 > 0 so that

= 0 is an unstable xed point, and f

(1) = 1 < 0 so that x

) is said to be a xed point of (7.2) if f(x

) = 0. Again, the local stability of a xed point can be determined by

+ (t) and y(t) = y

+ (t), where and

). Neglecting higher-order terms in the Taylor series, we

are = 3, 2 so that the xed point (0, 0) is an unstable

are = 1, 2 so that the xed point (0, 2) is a stable

are = 3, 1 so that the xed point (3, 0) is also a

2 can determined from the

2/2 and the incoming trajectories have positive slope

r. Clearly, two real xed points exist when

= r. The derivative of the right-hand-side is

(x) = r 2x, so that f

(r) = r. Therefore, for r < 0, x

= r is unstable, while for r > 0, x

r) = 2r. Therefore, the xed

r, the latter xed points existing

r) = 2r. Therefore, the xed point x

r exist and are

= 0 exists for all r, and four additional xed points can be

is real. Clearly, for the inner square-root to

1 + 4r becomes negative for r > 0.

= 0 (one xed point);

1 + 4r) (ve xed points);

1 + 4r) (three xed points).

) < 0. The minus root exists

) > 0. We summarize the

= 0. However, a catastrophe occurs as as soon as r > 0.

= 0, valid for all c, and the solutions to

> 0 only when (1 ) >